HK Malik

Engineering
Physics
Second Edition
About the Authors
Hitendra K Malik is currently Professor of Physics at the Indian Institute of Technology Delhi,
from where he received his PhD degree in the field of Plasma Physics, in 1995 at the age of 24. He
has been a merit scholarship holder throughout his academic career. He is the recipient of Career
Award from AICTE, Government of India, for his teaching and research, Outstanding Scientist
Award from VIF, India, for his contributions to Science, and 2017 Albert Nelson Marquis Lifetime
Achievement Award from USA. In addition, he received the prestigious Erasmus Mundus Visiting
Fellowship from European Union (Germany and France), JSPS Fellowship (two times) from Japan,
FRD Fellowship from South Africa and DAAD Fellowship from Germany. Owing to his worldwide recognition, his
name has been included in ‘Marquis Who’s Who’ in 2011, published from USA. Based on the survey conducted by
ResearchGate (RG), his scientific score has been found within top 5% of the scientists and researchers all over the world.
Professor Malik is highly cited in India and abroad for his research work and books with h-index of 24 and
i10-index of 70. Governments of India, Germany and France, through DST, CSIR, DRDO, AICTE, DAAD,
CEFIPRA, etc., have provided him funding to accomplish 12 sponsored research projects. He is on the editorial
board of 5 reputed research journals (including Springer). In recognition of his outstanding research and teaching
contributions, he has been asked to deliver more than 50 keynote and invited talks in India, Japan, South Korea,
USA, France, Germany, South Africa, and Turkey. Also, he has been chief guest in various universities, mentor of
faculty colleagues of engineering institutions, and member of organizing and advisory committees of national and
international conferences held in India and abroad.
He has guided 80 PhD, postgraduate and undergraduate theses, including 22 PhD theses in the area of laser/
microwave plasma interactions, particle acceleration, solitons, Terahertz radiation, Hall thrusters, plasma material
interaction, and nanotechnology. He has published more than 330 scientific papers in high impact factor journals
and conferences, including 19 independent articles. He has been reviewer for 72 Journals of international repute,
several sponsored research projects (Indian and Foreign agencies), and 18 PhD theses. He is an expert member of
academic and administrative bodies of 14 different universities and institutions from 8 states of India including UGC.
Apart from this book, he has also authored another textbook on Laser-Matter Interaction, CRC Press, 3 Chapters
in the Books Wave Propagation, InTechOpen Science, Croatia (featured as highly downloaded chapter), Society,
Sustainability and Environment, Shivalik Prakashan, New Delhi, and Plasma Science and Nanotechnology, Apple
Academic Press, exclusive worldwide distribution by CRC Press, a Taylor & Francis Group.
Ajay Kumar Singh has almost two decades of teaching experience in several engineering
institutions across North India. Currently, he is Professor of Physics at the Department of
Applied Sciences, Maharaja Surajmal Institute of Technology (MSIT), Janakpuri, New Delhi.
He has also served as the Head of Department at MSIT. Earlier, he was associated as Professor
(2003–2012) at the Department of Applied Science and Humanities, Dronacharya College of
Engineering, Haryana.
Dr. Singh completed his PhD from Aligarh Muslim University in the year 1999. During his
PhD, his work specifically focused on Uranium concentration in rock samples, soil samples
and fly ash samples. He also investigated radon levels in low and high background areas. He has published more
than 20 research papers and several articles in national and international journals and conferences. He has edited
and co-authored several books on Environment, Water Resources, Nuclear Physics, and Engineering Physics. His
book on Engineering Physics Practical and Tutorials has been highly appreciated by students. He is a life member
of Plasma Science Society of India (PSSI).
Dr. Singh was the ‘B.Tech First Year Syllabus Revision Committee’ coordinator representing all the affiliated
engineering colleges of Guru Gobind Singh Indraprastha University (GGSIPU). He is academic coordinator of PhD
scholars enrolled at MSIT under University School of Information, Communication and Technology, GGSIPU,
Dwarka, New Delhi. He is also supervisor of PhD students under USICT, which is a premier constituent institute
of GGSIPU. Dr. Singh has also been teaching a special course on Nanotechnology for the USICT PhD students. He
is also the teacher representative in the governing board of Maharaja Surajmal Institute of Technology.
Engineering
Physics
Second Edition
Hitendra K MaliK
Professor, Department of Physics
Indian Institute of Technology
Delhi
ajay KuMar SingH
Professor, Department of Applied Sciences
Maharaja Surajmal Institute of Technology
(MSIT) New Delhi
McGraw Hill Education (India) Private Limited

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Dedicated to
OMENDRA Bhaiya and all those moments
that remain with me as a source of inspiration
and help me to move ahead with great success,
satisfaction and optimistic approach
Brief Contents
Foreword xix
Preface to the Second Edition xxi
Preface to the First Edition xxiii
1. Interference 1
2. Diffraction 63
3. Polarisation 121
4. Lasers and Holography 155
5. Fibre Optics 186
6. Electron Optics 208
7. Waves and Oscillations 233
8. Simple Harmonic Motion and Sound Waves 259
9. Sound Waves and Acoustics of Buildings 284
10. Dielectrics 313
11. Electromagnetism 328
12. Theory of Relativity 395
13. Applied Nuclear Physics 451
14. Crystal Structure 517
15. Development of Quantum Mechanics 551
16. Quantum Mechanics 595
17. Free Electron Theory 634
18. Band Theory of Solids and Photoconductivity 654
19. Magnetic Properties of Solids 685
20. Superconductivity 716
21. X-Rays 735
viii Brief Contents
22. Nanoscience and Nanotechnology 751

Appendices
Appendix 1 Measurements and Errors 779
Appendix 2 Optics 790
Appendix 3 Mechanical Properties of Materials 794
Chapter-wise Answers to Objective Type Questions 805
Index 809
Contents ix
Contents
Foreword xix
Preface to the First Edition xxiii
1. Interference 1
Learning Objectives 1
1.1 Young’s Double Slit Experiment 2
1.2 Concept of Waves and Huygens’ Principle 2
1.3 Phase Difference and Path Difference 3
1.4 Coherence 4
1.5 Coherent Sources 5
1.6 Analytical Treatment of Interference 6
1.7 Conditions for Sustained Interference 8
1.8 Multiple Beam Superposition 9
1.9 Interference by Division of Wavefront 10
1.10 Interference by Division of Amplitude 16
1.11 Applications of Interference in the Field of Engineering 32
1.12 Scientific Applications of Interference 35
1.13 Homodyne and Heterodyne Detection 35
Summary 37
Solved Examples 38
Objective Type Questions 57
Short-Answer Questions 58
Practice Problems 58
Unsolved Questions 60
2. Diffraction 63
2.1 Young’s Double Slit Experiment: Diffraction or Interference? 64
2.2 Difference between Diffraction and Interference 64
2.3 Types of Diffraction 64
2.4 Fresnel’s Half-period Zones 66
2.5 Zone Plate 68
2.6 Fresnel’s Diffraction by a Circular Aperture 72
x Contents
2.7 Fraunhofer Diffraction by a Single Slit 76

2.8 Fraunhofer Diffraction by Double Slits 80
2.9 Fraunhofer Diffraction by N Slits: Diffraction Grating 84
2.10 Application of Diffraction Grating 90
2.11 Resolving Power of an Optical Instrument 91
2.12 Resolving Power of a Telescope 93
2.13 Resolving Power of a Microscope 94
2.14 Resolving Power of a Plane Diffraction Grating 96
2.15 Dispersive Power of a Plane Diffraction Grating 97
Summary 98
Solved Examples 99
3. Polarisation 121
3.1 Mechanical Experiment Showing Polarisation of Transverse Wave 122
3.2 Difference between Unpolarised Light and Polarised Light 123
3.3 Means of Production of Plane–Polarised Light 123
3.4 Theory of Production of Plane, Circularly and Elliptically Polarised Light 135
3.5 Optical Activity 138
3.6 Specific Rotation 138
3.7 Laurent’s Half-shade Polarimeter 138
3.8 Biquartz Polarimeter 139
3.9 Saccharimeter 140
3.10 Photoelasticity 141
Summary 142
Solved Examples 143
4. Lasers and Holography 155
4.1 Absorption and Emission of Radiation 156
4.2 Population Inversion 159
4.3 Characteristic of Laser Light 161
4.4 Main Components of Laser 163
4.5 Types of Laser 163
4.6 Applications of Lasers 169
4.7 Laser Cooling 170
4.8 Holography 170
Contents xi
4.9 Holography Versus Conventional Photography 171

4.10 Recording and Reconstruction of Image on Holograph 172
4.11 Types of Holograms 174
4.12 Applications of Holography 177
Summary 178
Solved Examples 179
5. Fibre Optics 186
5.1 Fundamental Ideas about Optical Fibre 187
5.2 Optical Fibres as a Dielectric Waveguide 188
5.3 Types of Optical Fibres 188
5.4 Acceptance Angle and Numerical Aperture 190
5.5 Fibre Optics Communication 193
5.6 Optical Fibre Sensors 196
5.7 Optical Fibre Connector 197
5.8 Optical Fibre Couplers 197
5.9 Applications of Optical Fibre Couplers 198
Summary 198
Solved Examples 200
6. Electron Optics 208
6.1 Specific Charge of an Electron 208
6.2 Determination of Specific Charge of an Electron: Thomson’s Method 209
6.3 Motion of an Electron in Uniform Electric and Magnetic Fields 210
6.4 Electrostatic and Magnetostatic Focusing 219
6.5 Scanning Electron Microscope (SEM) 220
6.6 Scanning Tunneling Microscope (STM) 225
Summary 226
Solved Examples 227
7. Waves and Oscillations 233
7.1 Translational Motion 233
7.2 Vibrational or Oscillatory Motion 233
7.3 Simple Harmonic Motion (SHM) 234
7.4 Differential Equation of SHM and Its Solution 234
7.5 Simple Pendulum 239
xii Contents
7.6 Mass-string System 241

7.7 Damped Harmonic Oscillator 243
7.8 Attenuation Coefficients of a Vibrating System 245
7.9 Forced Vibrations 247
7.10 Resonance 250
Summary 251
Solved Examples 252
8. Simple Harmonic Motion and Sound Waves 259
8.1 Superposition of Two SHMs 260
8.2 Sound Wave and its Velocity 260
8.3 Sound Displacement 261
8.4 Standing Waves 261
8.5 Standing Waves in Air Columns 262
8.6 Doppler Effect 264
8.7 Supersonic and Shock Waves 266
8.8 Derivation of Sound Speed 267
8.9 Intensity of Sound and Its Derivation 268
8.10 Sound-Intensity Level 269
8.11 Interference of Sound Waves in Time: Beats 270
8.12 Relation between Displacement and Pressure Amplitude 271
8.13 Lissajous Figures 272
8.14 Endoscopy 274
Summary 275
Solved Examples 276
9. Sound Waves and Acoustics of Buildings 284
9.1 Audible, Ultrasonic and Infrasonic Waves 284
9.2 Production of Ultrasonic Waves 285
9.3 Absorption and Dispersion of Ultrasonic Waves 290
9.4 Detection of Ultrasonic Waves 291
9.5 Applications of Ultrasonic Waves 292
9.6 Types of Acoustics 295
9.7 Acoustics of Buildings 296
9.8 Factors Affecting the Architectural Acoustics 300
Summary 301
Solved Examples 302
Contents xiii

True or False 311
10. Dielectrics 313
10.1 Dielectric Constant 313
10.2 Types of Dielectrics 314
10.3 Polarisation of Dielectrics 315
10.4 Types of Polarisation 317
10.5 Gauss’s Law in Dielectrics 318
10.6 Dielectric Loss 319
10.7 Clausius-Mosotti Equation 321
Summary 322
Solved Examples 323
11. Electromagnetism 328
11.1 Charge Density 329
11.2 Del Operator 330
11.3 Gradient 331
11.4 Divergence 331
11.5 Curl 332
11.6 Fundamental Theorem of Calculus 333
11.7 Fundamental Theorem for Gradient 333
11.8 Gauss’s or Green’s Theorem 334
11.9 Stokes’ Theorem 335
11.10 Electric Field and Electric Potential 335
11.11 Poisson’s and Laplace’s Equations 336
11.12 Capacitor 337
Æ
11.13 Magnetic Flux Density (B ) 338
Æ
11.14 Magnetic Field Strength (H ) 338
11.15 Ampere’s Circuital Law 338
11.16 Electrostatic Boundary Conditions 339
11.17 Scalar and Vector Potentials 341
11.18 Continuity Equation 342
11.19 Maxwell’s Equations: Differential Form 343
11.20 Maxwell’s Equations: Integral Form 348
11.21 Significance of Maxwell’s Equations 350
11.22 Maxwell’s Displacement Current and Correction in Ampere’s Law 351
11.23 Eletromagnetic (EM) Wave Propagation in Free Space 352
11.24 Transverse Nature of Electromagnetic Waves 355
xiv Contents
11.25 Maxwell’s Equations in Isotropic Dielectric Medium: EM Wave Propogation 355

11.26 Maxwell’s Equations in Conducting Medium: EM Wave Propagation and Skin
Depth 357
11.27 Electromagnetic Energy Density 360
11.28 Poynting Vector and Poynting Theorem 361
11.29 Wave Propagation in Bounded System: Waveguide 363
11.30 Coaxial Cable 364
Summary 365
Solved Examples 369
12. Theory of Relativity 395
12.1 Frame of Reference 395
12.2 Galilean Transformation 396
12.3 Michelson-Morley Experiment 398
12.4 Postulates of Special Theory of Relativity 401
12.5 Lorentz Transformation 401
12.6 Length Contraction 403
12.7 Time Dilation 405
12.8 Addition of Velocities 406
12.9 Variation of Mass with Velocity 407
12.10 Einstein’s Mass Energy Relation 409
Summary 411
Solved Examples 413
13. Applied Nuclear Physics 451
13.1 Basic Properties of Nucleus 452
13.2 Nuclear Forces 454
13.3 Binding Energy of Nucleus 455
13.4 Nuclear Stability 457
13.5 Nuclear Models 457
13.6 Nuclear Shell Model 457
13.7 Nuclear Liquid Drop Model 459
13.8 Radioactivity 460
13.9 Discovery of Neutron 469
13.10 Nuclear Reactions: Conservation Laws 471
13.11 Nuclear Fission 473
Contents xv
13.12 Nuclear Fusion 476

13.13 Controlled Fusion 477
13.14 Particle Accelerators 481
Summary 487
Solved Examples 492
14. Crystal Structure 517
14.1 Types of Solids 518
14.2 Unit Cell 518
14.3 Types of Crystals 519
14.4 Translation Vectors 520
14.5 Lattice Planes 521
14.6 Miller Indices 521
14.7 Simple Crystal Structures 523
14.8 Coordination Number 523
14.9 Interplanar Spacing 524
14.10 Nearest Neighbour Distance and Atomic Radius 525
14.11 Packing Fraction 526
14.12 Potential Energy Curve and Nature of Interatomic Forces 528
14.13 Different Types of Bonding Forces 529
14.14 Crystal Structure Analysis 531
14.15 Point Defects in Solids 534
Summary 537
Solved Examples 539
15. Development of Quantum Mechanics 551
15.1 Blackbody Radiation: Spectral Distribution 552
15.2 Planck’s Quantum Hypothesis 554
15.3 Simple Concept of Quantum Theory 556
15.4 Wave Particle Duality 557
15.5 Photoelectric Effect 557
15.6 de Broglie Waves: Matter Waves 559
15.7 Compton Effect: Compton Scattering 561
15.8 Phase and Group Velocities: de Broglie Waves 565
Summary 569
xvi Contents
Solved Examples 570

16. Quantum Mechanics 595
16.1 Heisenberg Uncertainty Principle 596
16.2 Wave Function and Its Physical Significance 601
16.3 Time Independent Schrödinger Equation 602
16.4 Time Dependent Schrödinger Equation 603
16.5 Operators 604
16.6 Applications of Schrödinger Equation 605
16.7 Quantum Statistics 614
Summary 617
Solved Examples 618
17. Free Electron Theory 634
17.1 Lorentz–Drude Theory: Classical Free Electron Theory of Metals 635
17.2 Applications of Lorentz–Drude Theory 636
17.3 Limitations of Lorentz–Drude or Free Electron Theory 638
17.4 Quantum Theory of Free Electrons 638
17.5 Thermionic Emission 644
Summary 645
Solved Examples 646
18. Band Theory of Solids and Photoconductivity 654
18.1 Kronig-Penney Model 655
18.2 One- and Two-Dimensional Brillouin Zones 659
18.3 Effective Mass of an Electron 660
18.4 Distinction between Insulators, Semiconductors and Conductors (Metals) 661
18.5 Intrinsic Semiconductor 664
18.6 Extrinsic Semiconductor 667
18.7 Hall Effect 669
18.8 Photoconductivity 671
18.9 Simple Model of Photoconductor 671
Contents xvii
18.10 Effect of Traps 674

18.11 Applications of Photoconductivity 675
Summary 676
Solved Examples 678
19. Magnetic Properties of Solids 685
19.1 Magnetic Moment of an Electron 687
19.2 Classification of Magnetic Materials 688
19.3 Comparison of Properties of Paramagnetic, Diamagnetic and Ferromagnetic
Materials 689
19.4 Classical Theory of Diamagnetism (Langevin’s Theory) 689
19.5 Classical Theory of Paramagnetism (Langevin’s Theory) 693
19.6 Classical Theory of Ferromagnetism 696
Æ Æ
19.7 Hysteresis: Nonlinear Relationship between B and H 697
19.8 Energy Loss Due to Hysteresis 698
19.9 Importance of Hysteresis Curve 699
19.10 Magnetic Circuits 700
19.11 Forces on Magnetic Materials 700
19.12 Magnetic Materials and Their Applications 700
Summary 702
Solved Examples 705
20. Superconductivity 716
20.1 Electrical Resistivity of Solids and Phonons 717
20.2 Properties of Superconductors 717
20.3 Classification of Superconductors 719
20.4 Effect of Magnetic Field 719
20.5 Isotope Effect 720
20.6 London Equations 720
20.7 Penetration Depth 722
20.8 Cooper Pairs 723
20.9 Bose–Einstein Condensation 724
20.10 BCS Theory: Qualitative Explanation 724
20.11 Coherence Length 725
20.12 High Temperature (Hi-Tc) Superconductivity 726
20.13 Application of Superconductivity 726
Summary 727
Solved Examples 729
xviii Contents

21. X-Rays 735
21.1 Origin of X-rays 736
21.2 Properties of X-rays 736
21.3 X-ray Spectra 737
21.4 Moseley’s Law 738
21.5 Practical Applications of X-rays 740
Summary 740
Solved Examples 741
22. Nanoscience and Nanotechnology 751
22.1 How Nanomaterials are Differenet from Bulk Materials? 752
22.2 Difference between Nanoscience and Nanotechnology 752
22.3 Quantum Confinement and Classification of Nanostructures 752
22.4 Nanoscale in 1-D 753
22.7 Applications of Nanotechnology 768
22.8 Limitations of Nanotechnology 770
22.9 Disadvantages of Nanotechnology 770
Summary 771
Solved Examples 773
Appendices
Appendix 1 Measurements and Errors 779
Appendix 2 Optics 790
Appendix 3 Mechanical Properties of Materials 794
Index 809
Foreword xix
Foreword
It gives me immense pleasure to see the present textbook on “Engineering Physics” which
covers almost the entire syllabus taught at undergraduate level at different engineering
colleges and institutions throughout India. I complement the authors and appreciate their
efforts in bringing out this book written in a very simple language. The text is comprehensive
and the explanation of topics is commendable. I understand that this book carries all the
elements required for a good presentation.
I have been a student of IIT Kharagpur and later on taught at IIT Delhi. Being a part of
the IIT system, I recognise that the rigorous and enriching teaching experience at IITs originating from the
interaction with the best engineering students and their strong feedback results in continuous evolution and
refinement of the teachers. This spirit is reflected in the comprehensive and in-depth handling of important
topics in a very simple manner in this book. I am happy to note that this textbook has been penned down by
IITian and hope that it would serve to be a good textbook on the subject. Since this book also covers advanced
topics, it will be an important learning resource for the teachers, and those students who wish to develop
research skills and pursue higher studies. I hope that the book is well received in the academic world.
Professor Prem Vrat

Vice-Chancellor, U.P. Technical University, Lucknow
Founder Director, IIT Roorkee
Preface to the Second Edition
The first edition of the textbook was appreciated by the teachers and students of many universities, engineering
colleges and institutes, including IIT’s throughout India. Words of appreciation were also received from
faculty colleagues from Japan, China, Taiwan, Russia, Canada, South Korea, Pakistan, Bangladesh, Turkey,
Iran, South Africa, Germany, France, United Kingdom, and United States of America. Students preparing for
GATE/CSIR competitive examinations also suggested for more examples in the book and inclusion of topics
of postgraduate level. The students very enthusiastically informed us about the utility of the book for the
preparation of interviews for admission in PhD programmes at IITs and other universities (including foreign
universities) or to get government jobs in India.
In view of all the above points, we have come up with the second edition of the book, where we have used
simple language for explaining each and every topic. We have included more physical insight, wherever
required. Some chapters are thoroughly revised in terms of new topics and solved problems. We have also
updated advanced topics keeping in mind the research going on in these fields. The solutions to the Objective-
Type Questions are also provided at the end of the book.
In particular, Chapter 4 includes details of the topic Population Inversion which covers various schemes
for the same, i.e., two-level, three-level and four-level systems. In Chapter 5, a topic on Optical Fibres as a
Dielectric Waveguide is included. After Chapter 7 on Waves and Oscillations, a new Chapter 8 on Simple
Harmonic Motion and Sound Waves has been included that discusses standing waves, supersonic and shock
waves, in addition to sound waves, Doppler effect and Lissajous figures. Chapter 9 on Sound Waves and
Acoustics of Buildings has been thoroughly revised. In this chapter, Recording and Reproduction of Sound
has been withdrawn and other topics are revisited. New topics on ultrasonics have been included which talk
about production of ultrasonic waves and their absorption, dispersion, detection and applications. In Chapter
10 on Dielectrics, a topic Energy Stored in an Electrostatic Field is withdrawn as its concept is discussed
in Chapter 11 on Electromagnetism. Moreover, details of Clausius-Mosotti equation are revised with the
inclusion of physical insight of this equation. The chapter on Electromagnetism has been thoroughly revised.
For example, Section 11.21 has been rewritten in order to make the readers understand which form of the
Maxwell’s equations is appropriate for free space, dielectric medium and conducting medium and how are
these equations modified in these media. Bound charges and bound currents are also discussed. The solution to
wave equation in conducting medium is included as Section 11.28.1, where dispersion relation, skin depth and
phase relationship of the electric and magnetic field vectors are discussed. New solved problems, objective-
type questions and other practice problems are also included in order to provide an indepth knowledge on the
electromagnetic fields and their propagation in different media.
In Chapter 12 on Theory of Relativity, physical insight to two interesting topics, viz. Length Contraction
and Time Dilation is provided. Several new solved problems on various topics are also provided for the
readers. Chapter 13 on Applied Nuclear Physics has been thoroughly revised and new topics are included on
xxii Preface to the Second Edition
basic properties of nucleus, nuclear forces, binding energy of nucleus, nuclear stability and various nuclear
models, in addition to more equations and problems, both solved and unsolved. Introduction part of Chapter
16 on Quantum Mechanics has been revised. The topic on Thermionic Emission (Section 17.7) has been
shortened but significance of Richardson’s equation is included. The earlier Chapter 21 on Photoconductivity
and Photovoltaics has been withdrawn but its important topics, viz. photoconductivity, simple model of
photoconductor and effect of traps, are included in Chapter 18 on Bond Theory of Solids and Photoconductivity.
The much important Chapter 22 on Nanophysics has been rewritten in view of recent advances in the
field. Now, it is renamed as Nanoscience and Nanotechnology. Certain new topics are included to clarify
how nanomaterials are different from bulk materials and to know the differences between nanoscience and
nanotechnology. The chapter very systematically discusses the nanoscales in 1D, 2D, 3D and OD. Particu-
larly, nanowires, carbon nanotubes, inorganic nanotubes, biopolymers, nanoparticles, buckyballs/fullerenes
and quantum dots are discussed in detail along with the methods of their synthesis, properties and their
applications. Finally, the applications, limitations and disadvantages of nanotechnology are also discussed.
The exhaustive OLC supplements of the book can be accessed at http://www.mhhe.com/malik/ep and contain
the following:
For Instructors
• Solution Manual
• Chapter-wise Power Point slides with diagrams and notes for effective lecture presentations
For Students
• A sample chapter
• A Solved Question Paper
• An e-guide to aid last minute revision need
We believe the readers shall find the second edition of the book more beneficial in terms of syllabus covered,
quality of topics, large number of solved problems aimed at providing physical insight to various topics,
and teaching various methods of solving difficult problems. The systematic approach adopted in the present
book shall certainly help the teachers and students providing for crystal clear understanding of the topics and
carrying out research in the related fields. This edition will be vital in enhancing the self confidence of our
UG and PG students which will help them in advancing their careers.
Finally, we look forward to receive feedback from the teachers and students on the recent edition of the book.
H K Malik
Ajay K Singh
Publisher’s Note:
McGraw Hill Education (India) invites suggestions and comments, all of which can be sent to
info.india@mheducation.com (kindly mention the title and author name in the subject line).
Piracy-related issues may also be reported.
Preface to the First Edition
Physics is a mandatory subject for all engineering students, where almost all the important elements of
the subject are covered. Finally, these evolve as different branches of the engineering course. The book
entitled Engineering Physics has been written keeping in mind the need of undergraduate students from
various engineering and science colleges of all Indian universities. It caters to the complete syllabus for
both–Physics-I and Physics-II papers in the first year Engineering Physics course.
The aim of writing this book has been to present the material in a concise and very simple way so that even
weak students can grasp the fundamentals. In view of this, every chapter starts with a simple introduction
and then related topics are covered with a detailed description along with the help of figures. Particularly the
solved problems (compiled from University Question Papers) are at the end of each chapter. These problems
are not merely numerical; many of them focus on reasoning and require thoughtful analysis. Finally, the chap-
ters carry unsolved questions based on which the students would be able to test their knowledge as to what
they have acquired after going through various chapters. A chapter-end summary and list of important formu-
lae will be helpful to students for a quick review during examinations. The rich pedagogy consists of solved
examples (450), objective-type questions (230), short-answer questions (224) and practice problems (617).
The manuscript has been formulated in such a way that students shall grasp the subject easily and save their
time as well. Since the complete syllabus is covered in a single book, it would be highly convenient to both.
The manuscript contains 22 chapters which have been prepared as per the syllabus taught in various colleges
and institutions. In particular, the manuscript discusses optics, lasers, holography, fibre optics, waves,
acoustics of buildings, electromagnetism, theory of relativity, nuclear physics, solid state physics, quantum
physics, magnetic properties of solids, superconductivity, photoconductivity and photovoltaic, X-rays and
nanophysics in a systematic manner. We have discussed advanced topics such as laser cooling, Bose-Einstein
condensation, scanning electron microscope (SEM), scanning tunnelling microscope (STM), controlled
fusion including plasma, Lawson criterion, inertial confinement fusion (ICF), plasma based accelerators,
namely, plasma wake field accelerator, plasma beat wave accelerator, laser wake field accelerator and self-
modulated laser wake field accelerator, and nanophysics with special emphasis on properties of nanoparticles,
carbon nanotubes, synthesis of nanoparticles and applications of nanotechnology. These will be of interest to
the teachers who are involved in teaching postgraduate courses at the universities and the students who opt for
higher studies and research as their career. Moreover, a series of review questions and problems at the end of
each chapter together with the solved questions would serve as a question bank for the students preparing for
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xxiv Preface to the First Edition
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Shyam Singh Aryavart Institute of Technology and Management, Lucknow
R S Tiwari Apollo Institute of Engineering, Kanpur
Kamlesh Pathak SVNIT, Surat, Gujarat
Kanti Jotania M S University, Baroda, Gujarat
Vijayalakshmi Sanyal Bharathiyar College of Engineering and Technology, Karaikal, Tamil Nadu
A K Meikap NIT, Durgapur, West Bengal
K Sivakumar Anna University, Chennai
H K Malik
Ajay K Singh
Interference 1
Learning Objectives
After reading this chapter you will be able to LO5 Discuss analytical treatment of
interference and conditions for
LO1 Explain interference through Young’s sustained interference
double slit experiment
LO6 Examine multiple beam superposition
LO2 Describe the concept of wave and and interference by division of
Huygen’s principle wavefront and amplitude
LO3 Illustrate phase and path difference LO7 Review engineering/scientific
LO4 Explain coherence, its various types applications of interferences including
and coherent sources homodyne and heterodyne detection
Introduction
You would have seen beautiful colours in soap films or patch of oil floating on the surface of water.
Moreover, the colour gets changed when you watch it from different angles. Did you ever try to find
out the reason? In scientific language, this takes place due to the phenomenon of interference. The
phenomenon of interference of light tells us about the wave nature of the light. In optics, the interference
means the superposition of two or more waves which results in a new wave pattern. Here, we are talking
about the interaction of waves emerging from the same source or when the frequencies of these waves
are the same. In the context of light, which is an electromagnetic wave, we say that when the light from
two different sources moves in the same direction, then these light wave trains superimpose upon each
other. This results in the modification of distribution of intensity of light. According to the principle of
superposition, this is called the interference of light. More precisely the interference can be defined as the
interaction between two or more waves of the same or very close frequencies emitted from coherent
sources (defined later), where the wavefronts are combined according to the principle of superposition.
The resulting variation in the disturbances produced by the waves is called the interference pattern.
Thomas Young, in 1802, explained the interference successfully in his double slit experiment.
2 Engineering Physics
1.1 YOung’s DOubLe sLit experiment LO1

The phenomenon of interference may be better understood by taking two point light sources S1 and S2 which
produce similar waves (Fig. 1.1). Let the sources S1 and S2 be at equal distances from the main source S while
being close to each other. Since the sources emit waves in all the directions, the spherical waves first pass
through S and then S1 and S2. Finally these waves expand into the space.
The crests of the waves are represented by complete arcs and the troughs
by dotted arcs. It is seen that constructive interference takes place at the
points where the crests due to one source meet the crests due to another
source or where their troughs meet each other. In this case, the resultant
S1
amplitude will be the sum of the amplitudes of the separate waves and
hence the intensity of the light will be maximum at these points. Similarly, S
at those points where crests due to one source meet the troughs due to S2
another source or vice-versa, the resultant amplitude will be the difference
of the amplitudes of the separate waves. At these points the intensity of
the waves (or light) will be minimum. Therefore, due to the intersection
of these lines, an alternate bright and dark fringes are observed on the Screen
screen placed at the right side of the sources S1 and S2. These fringes are
Figure 1.1
obtained due to the phenomenon of interference of light.
1.2 COnCept Of Waves anD HuYgens’ prinCipLe LO2

A wave is a disturbance that propagates through space and time, usually with the transference of energy from
one point to another without any particle of the medium being permanently displaced. Under this situation,
the particles only oscillate about their equilibrium positions. If the oscillations of the particles are in the
direction of wave propagation, then the wave is called longitudinal wave. However, if these oscillations take
place in perpendicular direction with the direction of wave propagation, the wave is said to be transverse
in nature. In electromagnetic waves, such as light waves, it is the changes in electric and magnetic fields
which represent the wave disturbance. The progress of the wave propagation is described by the passage of a
waveform through the medium with a certain velocity called the phase velocity or wave velocity. However,
the energy is transferred at the group velocity of the waves making the waveform.
The wave theory of the light was proposed in 1678 by Huygens, a Dutch scientist. On the basis of his wave
theory, he explained satisfactorily the phenomena of reflections, refraction etc. In the beginning, Huygens’
supposed that these waves are longitudinal waves but later he came to know that these waves are transverse
in nature. Huygens’ gave a hypothesis for geometrical construction of the position of a common wavefront
at any instant when the propagation of waves takes place in a medium. The wavefront is an imaginary
surface joining the points of constant phase in a wave propagated through the medium. The way in which the
wavefront is propagated further in the medium is given by Huygens’ principle. This principle is based on the
following assumptions:
(i) Each point on the given wavefront acts as a source of secondary wavelets.
(ii) The secondary wavelets from each point travel through space in all the directions with velocity of
light.
(iii) A surface touching the secondary wavelets tangentially in the forward direction at any given time
constructs the new wavefront at that instant. This is known as secondary wavefront.
Interference 3
In order to demonstrate the Huygens’ principle, we consider P1 P2 P P1

the propagation of a spherical wavefront (Fig. 1.2a) or plane P
wavefront (Fig. 1.2b) in an isotropic (uniform) medium P2
1 1
(for example, ether) emerging from a source of light S. At
any time, suppose PQ is a section of the primary wavefront 2 2
drawn from the source S. To find the position of the wavefront
after an interval t, we take points 1, 2, 3, ... on the primary 3 3
wavefront PQ. As per Huygens’ principle, these points act S
4
4
as the source of secondary wavelets. Taking each point as
the centre, we draw spheres of radii ct, where c is the speed 5 5
of light. These spherical surfaces represent the position of 6
6
secondary wavelets at time t. Further, we draw a surface
P1Q1 that touches tangentially all these secondary wavelets Q2
in the forward direction. This surface P1Q1 is the secondary Q Q2 Q Q1
Q1
wavefront. Another surface P2Q2 in the backward direction is
not called the secondary wavefront as there is no backward (a) (b)
flow of the energy during the propagation of the light waves. Figure 1.2
1.3 pHase DifferenCe anD patH DifferenCe LO3

As mentioned, the interference pattern is obtained when the two or more waves superimpose each other. In
order to understand this pattern it is very important to know about the path and phase differences between
the interfering waves.
1.3.1 phase Difference
Two waves that have the same frequencies and different phases are known to have a phase difference and are
said to be out of phase, with each other. If the phase difference is 180°, then the two waves are said to be in
antiphase and if it is 0°, then they are in phase as shown in Fig. 1.3(a and b). If the two interfering waves meet
at a point where they are in antiphase, then the destructive interference occurs. However, if these two waves
meet at a point where they are in the same phase, then the constructive interference takes place.
(a) (b)
Figure 1.3
1.3.2 path Difference
In Fig. 1.4, while the two wave crests are traveling a different distance from their sources, they meet at a point
P in such a way that a crest meets a crest. For this particular location on the pattern, the difference in distance
traveled is known as path difference.
1.3.3 Relation between Path Difference and

phase Difference
It is clear from the positions of crests or troughs of the waves r1
that if the path difference between the two waves is equal to S1 P
the wavelength l, the corresponding phase difference is 2p
(360°). Suppose for a path difference of d the corresponding
phase difference is f. Then it is clear that r2
S2
2p
f= d d = r2 − r1
l
2p (Path difference)
Phase difference = ¥ Path difference (i)
l Figure 1.4
This can be made clearer with the help of Fig. 1.4, where two sources of waves S1 and S2 are shown. The
wavelength of these sources is l and the sources are in phase at S1 and S2. The frequencies of both the waves
are taken to be the same as f. Therefore, the angular frequency w = 2pf. They travel at the same speed and the
2p
propagation constant for them is k = . We can write the wave equations for both the waves at point P as
l
y1 = a cos(wt – kr1) for the wave emerging from source S1 and
y2 = a cos(wt – kr2) for the wave emerging from source S2
Here (wt – kr1) is the phase f1 and (wt – kr2) is the phase f2. Therefore, the phase difference between them is
f1 – f2, given by f1 – f2 = wt – kr1 – wt + kr2 = k(r2 – r1).
2p
Using Eq. (i) and k = , the path difference is obtained as
l
Path difference d = r2 – r1.
1.4 COHerenCe LO4

Coherence is a property of waves that helps in getting stationary interference, i.e., the interference which is
temporally and spatially constant. During interference the waves add constructively or subtract destructively,
depending on their relative phases. Two waves are said to be coherent if they have a constant relative phase.
This also means that they have the same frequency. Actually the coherence is a measure of the correlation
that exists between the phases of the wave measured at different points. The coherence of a wave depends on
the characteristics of its source.
1.4.1 Temporal Coherence
Temporal coherence is a measure of the correlation between the phases of a wave (light) at different points
along the direction of wave propagation. If the phase difference of the wave crossing the two points lying
along the direction of wave propagation is independent of time, then the wave is said to have temporal
coherence. Temporal coherence is also known as longitudinal coherence. This tells us how monochromatic
a source is. In Fig. 1.5A, a wave traveling along the positive x-direction is shown, where two points A and B
are lying on the x-axis. Let the phases of the wave at these points at any instant t be fA and fB, respectively,
and at a later time t¢ they be f¢A and f¢B . Under this situation, if the phase difference fB – fA = f B¢ - f A¢ , then
the wave is said to have temporal coherence.
Interference 5
x axis
A B
Figure 1.5A
1.4.2 Spatial Coherence
Spatial coherence is a measure of the correlation between the phases of a wave (light) at different points
transverse to the direction of propagation. If the phase difference of the waves crossing the two points lying
on a plane perpendicular to the direction of wave propagation is independent of time, then the wave is said to
have spatial coherence. This tells us how uniform the phase of the wavefront is. In Fig. 1.5B, a wave traveling
along the positive x-direction is shown, where PQRS is a transverse plane and A and B are the two points
situated on this plane within the waveforms. Let the waves crossing these points at any time t have the same
phase f and at a later time t¢ the phases of the waves are again the same but equal to f¢. Under this situation,
the waves are said to have spatial coherence.
Q
P
A
x axis
B
x axis
R
S
Figure 1.5B
1.4.3 Coherence Time and Coherence Length

A monochromatic source of light emits radiation of a single frequency (or wavelength). In practice, however,
even the best source of light emits radiations with a finite range of wavelengths. For a single frequency wave,
the time interval over which the phase remains constant is called the coherence time. The coherence time is
generally represented by Dt. In a monochromatic sinusoidal wave the coherence time is infinity because the
phase remains constant throughout. However, practically the coherence time exists and the distance traveled
by the light pulses during this coherence time is known as coherence length DL. The coherence length is
also called the spatial interval, which is the length over which the phase of the wave remains constant. The
coherence length and coherence time are related to each other according to the following formula
DL = cDt
1.5 COHerent sOurCes LO4

Two sources of light are said to be coherent, if they emit waves of the same frequency (or wavelength), nearly
the same amplitude and maintain a constant phase difference between them. Laser is a good example of
coherent source. In actual practice, it is not possible to have two independent sources which are coherent. This
can be explained as follows. A source of light consists of large number of atoms. According to the atomic
theory, each atom consists of a central nucleus and the electrons revolve around the nucleus in different orbits.
When an atom gets sufficient energy by any means, its electrons jump from lower energy level to higher
energy level. This state of an atom is called an
excited state. The electron lives in this state only
for about 10–8 seconds. After this interval of time
the electrons fall back to the inner orbits. During
this process, the atoms radiate energy in the form of
light. Out of the large number of atoms some of
them emit light at any instant of time and at the next
instant other atoms do so and so on. This results in
the emission of light waves with different phases.
So, it is obvious that it is difficult to get coherent Many Source Points Many Wavelengths
light from different parts of the same source (Fig.
1.6). Therefore, two independent sources of light Figure 1.6
can never act as coherent sources.
1.5.1 Production of Coherent Light from Incoherent Sources

An ordinary light bulb is an example of an incoherent source. We can produce coherent light from such an
incoherent source, though we will have to a lot of the light. If we use spatially filter the light coming from
such source, we can increase the spatial coherence (Fig. 1.7). Further, spectrally filtering of the light increases
the temporal coherence. This way we can produce the coherent light from the incoherent source.
Spatial Filter Spectral Filter
Coherent
Light
Incoherent Source
Pinhole Wavelength Filter

Figure 1.7
1.6 anaLYtiCaL treatment Of interferenCe LO5

Let us consider the superposition of two waves of same frequency w and a constant phase difference f
traveling in the same direction. Their amplitudes are taken as a1 and a2, respectively. The displacement due
to one wave at any instant is given by
y1 = a1 sin wt (i)
and the displacement due to another wave at the same instant is given by
y2 = a2 sin (wt + f) (ii)
Interference 7
According to the principle of superposition, the resultant displacement (yR) is given by

yR = y1 + y2 (iii)
= a1 sin wt + a2 sin (wt + f)
= a1 sin wt + a2 sin wt cos f + a2 cos wt sin f
= (a1 + a2 cos f) sin wt + a2 sin f cos wt (iv)
Assuming a1 + a2 cos f = A cos q (v)
a2 sin f = A sin q (vi)
We obtain using Eq. (iv) – (vi)
yR = A sin (wt + q) (vii)
On squaring and adding Eqs. (v) and (vi), we have
2 2
A2 (sin2q+ cos2q) = a2 sin2f + a12 + 2a1a2 cos f + a2 cos2f
A2 = a12 + a22 (sin2f + cos2f) + 2a1a2 cos f (viii)
The resultant intensity is therefore given by
2 2
I = A2 = a1 + a2 + 2a1a2 cos f (ix)
The angle q can be calculated from Eqs. (v) and (vi) as
a2 sin f
tan q = (x)
a1 + a2 cos f
1.6.1 Condition for Constructive Interference

It is clear from Eq. (ix) that the intensity, I will be maximum at points where the values of cos f are +1, i.e,
phase difference f be 2np, with n = 0, 1, 2, 3, .... Then the maximum intensity is obtained from Eq. (ix) as
Imax = (a1 + a2)2 (xi)
In other words, the intensity will be maximum when the phase difference is an integral multiple of 2p. In this
case,
I max > a12 + a22
Thus, the resultant intensity will be greater than the sum of the individual intensities of the waves.
If a1 = a2 = a, then
Imax = 4a2
1.6.2 Condition for Destructive Interference

It is clear from Eq. (ix) that the intensity I will be minimum at points where cos f = –1. i.e., where phase
difference f = (2n + 1)p, with n = 0, 1, 2, 3,.... Then Eq. (ix) gives
Imin = (a1 – a2)2 (xii)
Therefore, it is clear that in destructive interference the intensity will be minimum when the phase difference
f is an odd multiple of p.
If a1 = a2, then Imin = 0
If a1 π a2, then Imin π 0
I min < a12 + a22

Thus, in the case of destructive interference the resultant intensity will be less than the sum of the individual
intensities of the waves.
Figure 1.8 represents the intensity variation with phase differences f graphically (for a1 = a2 = a).
I
4a2
–7p –6p –5p –4p –3p –2p –p 0 p 2p 3p 4p 5p 6p 7p

f
Figure 1.8
1.6.3 Conservation of Energy
The resultant intensity due to the interference of two waves a1 sin wt and a2 sin (wt + f) is given by Eq. (ix),
reproduced below
2 2
I = a1 + a2 + 2a1a2 cos f
2 2
\ Imax = a1 + a2 + 2a1a2 = (a1+ a2)2
and Imin = a12 + a22 – 2a1a2 = (a1 – a2)2
If a1 = a2 = a then
Imax = 4a2 and Imin = 0
Therefore, average intensity (Iav) will be obtained as

Iav = 2a2
2 2
For unequal amplitudes a1 and a2 the average intensity would be (a1 + a2 ) . Thus, in interference only some
part of energy is transferred from the position of minima to the position of maxima, and the average intensity
or energy remains constant. This shows that the phenomenon of interference is in accordance with the law of
conservation of energy.
1.7 COnDitiOns fOr sustaineD interferenCe LO5

Sustained interference means a constant interference of light waves. In order to obtain such interference, the
following conditions must be satisfied
(i) The two sources should emit waves of the same frequency (wavelength). If it is not so, then the
positions of maxima and minima will change with time.
Interference 9
(ii) The waves from the two sources should propagate along the same direction with equal speeds.
(iii) The phase difference between the two interfering waves should be zero or it should remain constant.
It means the sources emitting these waves must be coherent.
(iv) The two coherent sources should be very close to each other, otherwise the interference fringes will
be very close to each other due to the large path difference between the interfering waves. For the
large separation of the sources, the fringes may even overlap and the maxima and minima will not
appear distinctly.
(v) A reasonable distance between the sources and screen should be kept, as the maxima and minima
appear quite close if this distance is smaller. On the other hand, the large distance of the screen
reduces the intensity.
(vi) In order to obtain distinct and clear maxima and minima, the amplitudes of the two interfering waves
must be equal or nearly equal.
(vii) If the source is not narrow, it may act as a multi source. This will lead to a number of interference
patterns. Therefore, the coherent sources must be narrow.
(viii) In order to obtain the pattern with constant fringe width and good intensity fringes, the sources
should be monochromatic and the background should be dark.
1.7.1 Condition of Relative Phase Shift

This is regarding the introduction of additional phase change between the interfering waves when they emerge
after reflecting from two different surfaces. In most of the situations, the reflection takes place when the beam
propagates from the medium of lower refractive index to the medium of higher refractive index or vice-versa.
When the reflection occurs with light going from a lower index toward a higher index, the condition is called
internal reflection. However, when the reflection occurs for light going from a higher index toward a lower
index, the condition is referred to as external reflection. A relative phase shift of p takes place between the
externally and internally reflected beams so that an additional path difference of l/2 is introduced between
the two beams. If both the interfering beams get either internally or externally reflected, no phase shift takes
place between them.
1.8 muLtipLe beam superpOsitiOn LO6

In Section 1.6, we have given theoretical analysis of the interference due to the superposition of two waves
of the same frequency and the constant phase difference. The intensity of the interference pattern showed its
dependence on the amplitudes of the interfering waves. However, now we consider a large number of waves
of the same frequency and amplitude, which propagate in the same direction. The amount by which each
wave train is ahead or lags behind the other is a matter of chance. Based on the amplitude and intensity of
the resultant wave, we can examine the interference. We assume n number of wave trains whose individual
amplitudes are equal (= a, say). The amplitude of the resultant wave can be understood as the amplitude of
motion of a particle undergoing n simple harmonic motions (each of amplitude a) at once. In this case, if all
these motions are in the same phase, the resultant wave will have an amplitude equal to na and the intensity
would be n2a2, i.e., n2 times that of one wave. However, in our case, the phases are distributed purely at
random, as shown in Fig. 1.9 as per graphical method of compounding amplitudes. Here, the phases f1, f2,
f3, ... take arbitrary values between 0 and 2p. The intensity due to the superposition of such waves can be
calculated by the square of the resultant amplitude A. In order to find A2, we should square the sum of the
projections of all vectors a along the x-direction and add it to the square of the corresponding sum along the
y-direction. The summation of projections along x-direction are given by the following expression
a(cos f1 + cos f2 + cos f3 + ... + cos fn)
y
a f3
f2
A a
a
f1
O X
Figure 1.9
The square of quantity in the parentheses gives the terms of the form cos2 f1, 2 cos f1 cos f2, etc. It is seen that
the sum of these cross product terms increases approximately in proportion to number n. So we do not obtain
a definite result with one given array of arbitrarily distributed waves. For a large number of such arrays, we
find their average effect in computing the intensity in any physical problem. Under this situation, it is safe to
conclude that these cross product terms will average to zero. So we consider only the cos2 f terms. Similarly,
for the y projections of the vectors we obtain sin2 f terms. With this we have
I ª A2 = a2(cos2 f1 + cos2 f2 + cos2 f3 + ... + cos2 fn) + a2(sin2 f1 + sin2 f2 + sin2 f3 + ... + sin2 fn).
Using the identity sin2 fp + cos2 fp = 1, the above expression reduces to I ª a2 ¥ n.
Since a2 is the intensity due to a single wave, the above relation shows that the average intensity resulting
from the superposition of n waves with arbitrary phases is n times of a single wave. It means the resultant
amplitude A increases in proportion with in length as n gets increased.
1.9 interferenCe bY DivisiOn Of WavefrOnt LO6

This method uses multiple slits, lenses, prisms or mirrors for dividing a single wavefront laterally to form
two smaller segments that can interfere with each other. In the division of a wavefront, the interfering beams
of radiation that left the source in different directions and some optical means is used to bring the beams
back together. This method is useful with small sources. Double slit experiment is an excellent example of
interference by division of wavefront. Fresnel’s biprism is also used for getting interference pattern based on
this method.
1.9.1 Fresnel’s Biprism
Fresnel’s Biprism is a device by which we can obtain two virtual coherent sources of light to produce
sustained interference. It is the combination of two acute angled prisms which are joined with their bases in
such a way that one angle becomes obtuse angle q¢ of about 179° and remaining two angles are acute angles
each of about 1/2°, as shown in Fig. 1.10.
Interference 11
S1 d
d
B
q¢
2d
S
C
q
S2
D
A¢
Y
Figure 1.10
Let monochromatic light from slit S fall on the biprism, placed at a small distance from S. When the light falls
on upper part of the biprism, it bends downward and appears to come from source S1. Similarly, the other
part of the light when falls on the lower part of the biprism, bends upward and appears to come from source
S2. Here, the images S1 and S2 act as two virtual coherent sources of light (Fig. 1.10). Coherent sources are
the one that have a constant or zero phase difference throughout. In the situation, on placing the screen XY on
right side of the biprism, we obtain an alternate bright and dark fringes in the overlapping region BC.
1.9.1.1 Theory of Fringes
X
Let A and B be two virtual coherent sources of light separated by a distance
P
2d. The screen XY, on which the fringes are obtained, is separated by a
distance D from the two coherent sources, as shown in Fig. 1.11. The
point C on the screen is equidistant from A and B. Therefore, the path xn
difference between the two waves from sources A and B at point C is
zero. Thus the point C will be the centre of a bright fringe. On both sides A N
of C, alternately bright and dark fringes are produced. d
2d S C
Draw perpendiculars AN and BM from A and B on the screen. Let the d
B M
distance of a point P on the screen from the central bright fringe at C be D
xn. Y
From geometry, we have Figure 1.11
NP = xn – d; MP = xn + d
In right angled DANP,
AP2 = AN2 + NP2 (i)
2 2
= D + (xn – d)
È ( x - d )2 ˘
D 2 Í1 + n ˙
Î D2 ˚
1/2
È ( x - d )2 ˘
AP = D Í1 + n ˙
Î D2 ˚
È 1 ( xn - d ) 2 ˘
AP = D Í1 + ˙ , [as (xn - d ) << D ]
Î 2 D2 ˚
1 ( xn + d )2
AP = D + (ii) [By using Binomial Theorem]
2 D
Similarly, in DBMP,
1 ( xn + d ) 2 (iii)
BP = D +
2 D
Hence, the path difference between the waves reaching via AP and BP paths at the point P on the screen
È 2˘ È 2˘
D = BP - AP = Í D + 1 ( xn + d ) ˙ - Í D + 1 ( xn - d ) ˙
Î 2 D ˚ Î 2 D ˚
4x d
= n
2D
2d
D= xn (iv)
D
Condition for Bright Fringes: In order to interfere constructively and produce bright fringes, the two rays
should arrive at points P in phase. This is possible if the path difference is an integral multiple of l. Therefore,
D = nl
2d
xn = nl where n = 0, 1, 2 ……
D
nl D
xn = (v)
2d
Here it may be recalled that xn is the distance of the nth order bright fringe from the central maxima.
The distance of the next (n + 1)th maximum from the point C can be calculated by replacing n by n + 1 in
equation (v). Therefore,
lD
x( n +1) = (n + 1)
2d
The separation between two consecutive maxima gives the fringe width b, as follows
b = xn+1 – xn
or fringe width
lD
b= (vi)
2d
Interference 13
Condition for Dark Fringes: In order to interfere destructively and produce dark fringe at point P, the two
rays should arrive at this point in out of phase (phase difference of p). This is possible, if the path difference
l
is an odd multiple of . Therefore,
2
Ê 1ˆ
D = Á n + ˜ l , where n = 0, 1, 2, ...
Ë 2¯
From Eq. (iv)
2d l
D= xn = (2n + 1) (vii)
D 2
(2n + 1)l D
xn = (viii)
4d
Equation (viii) gives the distance of nth order dark fringe from the point C. The distance of the next (n+1)th
minimum from the point C will be
[2(n + 1) + 1]l D
x( n +1) =
4d
(2n + 3)l D
= (ix)
4d
Hence, the fringe width between two consecutive minima would be
(2n + 3)l D (2n + 1)l D
b = x( n +1) - xn = -
4d 4d
lD
b= (x)
2d
It is clear from Eqs. (vi) and (x) that the bright and dark fringes are of equal width.
1.9.1.2 Experimental Method for Determination of Wavelength of Light
The experimental setup used for the determination of wavelength of light consists of a good quality heavy
optical bench of about 1.5 meter length fitted with scale. It has four uprights that carry an adjustable slit S, a
biprism, a convex lens and a micrometer eyepiece, respectively. These components are shown in Fig. 1.12.
Each upright can be moved along the length of the optical bench and screws are provided to rotate the slit and
biprism in their own planes and the eyepiece can also move at right angle to the length of the optical bench.
To obtain well defined and sharp interference fringes, the following adjustments are necessary:
(i) Labeled optical bench by using spirit level and leveling screws.
(ii) Adjust all uprights to the same height.
(iii) Illuminate the vertical slit by monochromatic source of light. Make the slit narrow.
(iv) Now place the biprism on the second upright and try to adjust its edge parallel to the slit until two
equally bright virtual sources A and B are observed.
(v) Shift the micrometer eyepiece on the bench away from the slit and also move it at right angle to the
length of optical bench until the fringes are observed in the field of view.
(vi) In order to get fine fringes, change the position of the biprism slowly in its own plane such that its
edge remains parallel to the slit.
A
C1 E
S
C2
Figure 1.12
Lateral shift and its removal: On moving the micrometer eye piece on the bench towards the biprism, if
the fringes appear to shift at right angle to the optical bench then it is known as lateral shift (Fig. 1.13(a)).
However, if the principle axis and axis of optical bench become parallel, then no lateral shift remains, as
shown in Fig. 1.13(b).
Axis of Optical
Bench
Principle Axis
(a) Lateral shift (b) No lateral shift of fringes
Figure 1.13
1.9.1.3 Determination of Distance between Two Virtual Coherent Sources

For measuring 2d, a convex lens of short focal length is placed between the biprism and the micrometer eye
piece. This distance between the biprism and the micrometer eye piece is more than 4 times of the focal
length of the convex lens. By moving the lens we obtain two positions L1 and L2 of the convex lens such
that two separated images d1 and d2 of the two coherent sources respectively can be observed, as shown in
Fig. 1.14.
For the first position of lens, L1, the magnification is given as
v d1
= (i)
u 2d
and for second position of the lens, the magnification is
u d2
= (ii)
v 2d
Interference 15
Then from Eqs. (i) and (ii), we get

v u d1d 2
¥ = or (2d ) 2 = d1d 2
u v (2d ) 2
or 2d = d1d 2 (iii)
u v
2d d2 d1
L1 L2
v u
Figure 1.14
Therefore, the measurement of positions of images d1 and d2 will determine the distance 2d between the
sources. The wavelength l of monochromatic light can be calculated when we substitute the values of b, D
and 2d in the formula l = b(2d/D), derived in the previous section.
1.9.1.4 Determination of Thickness of Thin Transparent
Sheet (Displacement of Fringes)
Let A and B be two virtual coherent sources of light. The point C0 on the
P
screen is equidistant from both the sources (Fig. 1.15). When a transparent
material plate G of thickness t and having refractive index m, it is placed t
G
in the path of one of the light wave, we observe that the fringe which was xn
originally at C0 shifts to another position P, as shown in Fig. 1.15.
A
The time taken by the light wave from A to P partly through air and partly
through the plate is the same as the time taken by the other light wave C0
from B to P in air. If c and v be the velocity of light in air and in the plate,
B
respectively, then
Figure 1.15
BP AP - t t
= +
c c v
BP AP - t mt È c˘
or = + ÍÎQ m = v ˙˚
c c c
or BP = (AP – t) + mt
or BP – AP = (m – 1)t (i)
Here BP – AP is the path difference between the two interfering waves.
If the point P is originally occupied by the nth order bright fringe, then the path difference between the two
interfering waves will be
BP – AP = nl,
(m – 1) t = nl (ii)
The distance xn through which the fringe is shifted to point P from the central maximum C0 is given by
nl D
xn = (iii)
2d
lD
where, = b = fringe width.
2d
From Eq. (iii), we get
xn ◊ 2d
= nl (iv)
D
From Eqs. (ii) and (iv), we get
xn ◊ 2d
( m - 1)t =
D
xn ◊ 2d
or t= (v)
( m - 1) D
Therefore, by knowing xn, 2d, D and m, we can calculate thickness t of the glass plate by using Eq. (v).
1.10 interferenCe bY DivisiOn Of ampLituDe LO6

The method, which is used to produce two coherent sources from a common source, is called division of
amplitude that maintains the same width but reduced amplitude. After following different paths the two
waves of reduced amplitudes are combined to produce an interference pattern. In this method, the interfering
beams consist of radiation that has left the source in the same direction. This radiation is divided after leaving
the source and later combined to produce interference. This method can be used with extended sources.
Michelson interferometer is an example of interference by division of amplitude. Thin films are also used for
getting interference pattern based on this method.
1.10.1 Interference Due to Thin Films

This is clear that the interference takes place when the two waves superimpose each other after traveling
some distance, i.e., when there is a path difference between them. Since the thin film has its two surfaces, the
waves reflected from these surfaces can attain a path
difference and can interfere. The same may be applied C E
A M
on the waves that transmit through the film.
i i Air
1.10.1.1 Thin Film of Uniform Thickness (90 - i) D
t
B H
Consider a uniform transparent film having thickness
t and a refractive index m. A ray of light AB incident r r r
Film
at an angle i on the upper surface of the film is partly t
(µ > 1)
reflected along BC and partly refracted along BF at r r r r
an angle r. At point F the wave BF is again partly I G
reflected from the second surface along FD and partly F
t
i Air
emerges out along FK and so on. In this situation, the
interference occurs between reflected waves BC and J L
DE and also between the transmitted waves FK and K
GL (Fig. 1.16). Figure 1.16
Interference 17
The path difference between the reflected rays

D = (BF + FD)in film – (BM)in air
D = m (BF + FD) – BM
Q BF = FD
\ D = 2 mBF – BM (i)
In the right angled DBFH,
t t
cos r = or BF = (ii)
BF cos r
BH
and tan r = or BH = t tan r
t
BD = 2 ¥ BH
\ BD = 2t tan r (iii)
In the DBMD,
BM
sin i = or BM = BD sin i
BD
\ BM = 2t tan r sin i (iv)
From Eqs. (i), (ii) and (iv), we get
t
D = 2m - 2t tan r sin i (v)
cos r
sin i
Q m= or sin i = m sin r (vi)
sin r
2 mt sin r 2 mt
\ D= - 2t m sin r = [1 - sin 2 r ]
cos r cos r cos r
D = 2mt cos r (vii)
Equation (vii) represents only the apparent path difference and does not represent the effective total path
difference. When the light is reflected from the surface of an optically denser medium in case of rad BC, a
phase change of pa equivalent to path difference of l/2 is introduced. Therefore, the total path difference
between BC and DE will be
D = 2mt cos r + l/2 (viii)
Condition for Maxima: To have a maximum at a particular point, the two rays should arrive there in phase.
So the path difference must contain a whole number of wavelength, i.e.,
D = nl, n = 0, 1, 2..., (ix)
From Eq. (viii) and (ix), we get
2mt cos r + l/2 = nl
2mt cos r = nl – l/2
2mt cos r = (2n – 1)l/2 (x)
Condition for Minima: To have a minimum at a particular point, the two rays should arrive there in out of
phase (odd multiple of p) for which the path difference must contain a half odd integral number of wavelength,
i.e,
Ê 1ˆ
D = Án + ˜ l (xi)
Ë 2¯
Using Eq. (viii), we obtain
2mt cos r = nl where, n = 0, 1, 2, 3, … (xii)
It should be noted that the interference pattern will not be perfect because the intensities of the rays BC and
DE are not the same and their amplitudes are different.
In order to obtain the interference between the transmitted waves, we calculate the path difference between
the waves, FK and GL as under
D = (FD + DG)in film – (FJ)in air
D = m[FD + DG] – FJ
Q FD = DG
\ D = 2mFD – FJ (xiii)
DI t t
In DFDI, cos r = = or FD = (xiv)
FD FD cos r
FI FI
and tan r = = or FI = t tan r
DI t
FG = 2t tan r (xv)
In right angled DFJG,
FJ
sin i = or FJ = FG sin i
FG
\ FJ = 2t tan r sin i (xvi)
From Eq. (xiii), (xiv) and (xvi), we get
2mt
D= - 2t tan r sin i
cos r
2mt sin r È sin i ˘
=
cos r
- 2t
cos r
m sin r Í m = sin r ˙
Î ˚
2mt
= [1 - sin 2 r ] = 2 mt cos r
cos r
Since these two waves are emerging from the same medium, the additional phase difference (or path
difference) will not be introduced. Therefore, the total path difference
D = 2mt cos r (xvii)
Condition for Maxima: As discussed, it is possible when
D = nl (xviii)
Interference 19
From Eqs. (xvii) and (xviii), we get

2mt cos r = nl where, n = 0, 1, 2, 3, … (xix)
Condition for Minima: For obtaining minimum intensity, we should have
Ê 1ˆ
D = Án + ˜ l
Ë 2¯
Ê 1ˆ
which gives 2mt cos r = Á n + ˜ l where, n = 0, 1, 2, 3, …(xx)
Ë 2¯
Thus, the conditions for interference with transmitted light are obviously opposite to those obtained with
reflected light. Hence, if the film appears dark in the reflected light, it will appear bright in the transmitted
light and vice-versa. This shows that the interference pattern in the reflected and transmitted lights are
complimentary to each other.
(i) Necessity of an Extended Source of Light for Interference in Thin Films
When a thin transparent film is exposed to white light and seen in the reflected light, different colours are
seen in the film. These colours arise due to the interference of the light waves reflected from the top and
bottom surfaces of the film. The path difference between the reflected rays depends upon the thickness t,
refractive index m of the film and the angle q of inclination of the incident rays. The light which comes from
any point from the surface of the film will include the colour whose wavelength satisfies the equation 2mt cos
r = (2n – 1) l/2 and only this colour will be present with the maximum intensity in the reflected light.
When the transparent film of a large thickness as compared to the wavelength of the light, is illuminated by
white light, the path difference at any point of the film will be zero. In the case of such a thick film, at a given
point, the condition of constructive interference is satisfied by a large number of wavelengths, as << t. The
condition of destructive interference is also satisfied at the same point for the large number of wavelengths.
Therefore, consequently that point receives an average intensity due to the light of all wavelengths and no
colours are observed.
In the context of realization of above phenomena it is always needed to use a broad power of light that will
enable the eye to see whole of the film simultaneously.
If we use a point source, then we observe that different parts of reflected light cannot reach the eye due to
small size of the pupil, as shown in Fig. 1.17(a). The reflected rays only from a small portion of the film can
enter the eye. Hence, the whole of the film cannot be seen by the eye placed in a fixed position. However, if a
Eye Extended Source Eye
Point Source
(a) (b)
Figure 1.17
broad source of light is used to illuminate a thin film, the light reflected from each part of the film reaches the
eye placed in a fixed position, as shown in Fig. 1.17(b). Hence, one can seen the entire film simultaneously
by employing an extended source of light.
1.10.1.2 Non-uniform Thickness Film (Wedge Shaped Film)
Consider two plane surfaces OM and OM¢ inclined at an angle q enclosing a wedge shaped air film of
increasing thickness, as shown in Fig 1.18. A beam of monochromatic light is incident on the upper surface
of the film and the interference occurs between the rays reflected at its upper and lower surfaces. The
interference occurs between the reflected rays BK and DL, both of which are obtained from the same incident
ray of light AB.
K
E L
M¢
i D
r
A i (r + q)
i
B r+q t
N (r + q)
q r
O M
C H
q (r + q)
I¢
I
Figure 1.18
The path difference between the two reflected rays

D = [BC + CD]in film – [BE]in air
D = m(BC + CD) – BE
Q CD = CI
D = m(BC + CI) – BE
= mBI – BE
= m(BN + NI) – BE (i)
In right angled DBED,
BE
sin i = (ii)
BD
Interference 21
In right angled DBND,

BN
sin r = (iii)
BD
Dividing Eq. (ii) by Eq. (iii), we get
sin i BE sin i BE
= or =m=
sin r BN sin r BN
or BE = mBN (iv)
From Eqs. (i) and (iv), we get
D = m(BN + NI) – mBN
or D = mNI (v)
In right angled DDNI,
NI
cos (r + q ) =
DI
\ DI = DH + HI = t + t = 2t
NI
cos (r + q ) = or NI = 2t cos (r + q ) (vi)
DI
From Eqs. (v) and (vi), we get
D= 2mt cos (r + q) (vii)
Equation (vii), in the case of reflected light, does not represent the effective total path difference, as a phase
difference of p (Stokes phase change) has been introduced through the reflection of wave BK. Therefore, the
total path difference between the reflected rays,
D = 2mt cos(r + q) + l/2 (viii)
Equation (viii) shows that the path difference D depends on the thickness t. However, t is not uniform and it
is different at different positions.
At t = 0, Eq. (viii) reads
D = l/2
which is the condition for darkness. Therefore, the edge of the film appears to be dark. This is called zero
order band.
For normal incidence, i = 0 and r = 0. Then, the path difference
D = 2mt cos q + l/2 (ix)
Condition for Maxima: As explained earlier, the constructive interference takes place when
D = nl where, n = 0, 1, 2, 3, .... (x)
From Eqs. (ix) and (x), we get
2mt cos q + l/2 = nl
2mt cos q = (2n – 1) l/2 (xi)
Condition for Minima: In order to get destructive interference, the path difference
Ê 1ˆ
D = Án + ˜l (xii)
Ë 2¯
Ê 1ˆ
or 2 mt cos q + l /2 = Á n +˜l
Ë 2¯
2mt cos q = nl where, n = 0, 1, 2, 3, … (xiii)
(i) Nature of Fringes
For normal incidence of the light waves or a parallel incident beam, the incident angle remains constant and
hence the angle of refraction. If the light is monochromatic, then l is also fixed. Therefore, the change in path
difference will take place due to mt or thickness t of the film only. As we move outwards from the point of
contact O, the thickness of the film increases. However, at a particular place along a line parallel to the edge,
t has only one value. Since the locii of the points of constant thickness of the film are straight lines parallel to
the edge, straight bright and dark fringes parallel to the edge will be formed in the reflected light. If we use
the white light in place of monochromatic light, coloured fringes will be observed.
(ii) Derivation for Fringe Width
For a wedge shaped film the conditions of maxima and minima are reproduced below.
2mt cos (r + q) = (2n – 1)l/2
2mt cos (r + q) = nl
For normal incidence and small values of q the above conditions read as
2mt = (2n – 1)l/2 (xiv)
and 2mt = nl (xv)
If points A and C (Fig. 1.19) represent positions of two consecutive dark
fringes corresponding to film thicknesses AB = t1 and CD = t2 respectively, D
then the fringe width (w) will be equal to BE. Now from Eq. (xv), we get B q
E
the following condition corresponding to the points A and C.
q
2mt1 = nl and 2mt2 = (n + 1)l A C
or 2m(t2 – t1) = l Figure 1.19
or 2(CD – AB) = l
or 2(DE) = l (xvi)
But tan q = DE/BE or DE = BE tan q (xvii)
From Eqs. (xvi) and (xvii), we get
2m(BE tan q) = l
l
or BE = =w
2 m tan q
For smaller values of q, tan q ; q and we get
l
w= (xviii)
2 mq
Interference 23
It is clear from (xviii) that the fringe width w is independent of thickness t for smaller angle q. Therefore the
fringes are equally spaced and of same width for fixed l, m and q.
1.10.2 Newton’s Rings
If a plano-convex lens is placed such that its curved surface lies on a glass plate, then an air film of gradually
increasing thickness is formed between the two surfaces. When a beam of monochromatic (single wavelength)
light is allowed to fall normally on this film and viewed as shown in Fig. 1.20, an alternating dark and
bright circular fringes are observed. These circular fringes are formed because of the interference between
the reflected waves from the top and the bottom surfaces of the air film. These fringes are circular since
the air film has a circular symmetry and the thickness of the film corresponding to each fringe is same
throughout the circle. The interference fringes so formed were first investigated by Newton and hence known
as Newton’s rings.
The path difference between the two reflected rays, can be obtained as done in the case of wedge shaped film.
It is reproduced below as
D = 2mt cos (r + q) + l/2 (i)
Where (l/2) is due to Stokes phase change.
45°
Plano-Convex
Lens
Air Film
O
Glass Plate P
Figure 1.20
For normal incidence and an air film, i = 0, r = 0, m = 1. In addition, if q is also very small, then cos q = 1.
Under this situation, the path differences becomes
l
D = 2t + (ii)
2
Here t is the thickness of the air film at a particular point.
At the point of contact, t = 0
l
D=
2
which is the condition of minimum intensity and hence, the central spot of the ring will be dark.
Condition for Maxima: For constructive interference

D = nl (iii)
l
or 2t + = nl
2
l
or 2t = (2n – 1) where n = 0, 1, 2, 3... (iv)
2
Condition for Minima: For destructive interference
Ê 1ˆ
D = Án + ˜ l
Ë 2¯
l l
or 2t + = nl +
2 2
or 2t = nl where n = 0, 1, 2, 3… (v)
Diameter of Dark and Bright Rings: Let us consider the thickness of the air
film at point Q as t and r, as the radius of the fringe at that point together with R
as the radius of curvature of the lens (Fig. 1.21).
C
Hence, OC = CQ = R, HQ = rn
HC = R – t R
R
In right angled DCHQ H Q t
2 2 2
CQ = CH + HQ O P
2 rn
R2 = (R – t)2 + rn
or rn2 = 2Rt – t2 Figure 1.21
In actual practice, R is quite large and t is very small. Therefore, t2 may be neglected in comparison with 2Rt
\ rn2 = 2Rt
or rn2 = R ¥ 2t (vi)
For Bright Rings: From Eq. (iii), we get
l
2t = (2n - 1)
2
When we put this value of 2t in Eq. (vi), we get
l
rn2 = R ¥ (2n - 1)
2
2
Ê Dn ˆ l
ÁË ˜¯ = R ¥ (2n - 1)
2 2
or Dn2 = 2l R(2n - 1) (vii)
The above equation gives the diameter Dn of nth order bright fringe as
Dn = 2l R(2n - 1)
\ Dn µ (2n - 1) (viii)
Interference 25
Thus the diameter of the bright circular fringe(s) is proportional to the square root(s) of the odd natural
numbers.
For Dark Rings: Applying the condition 2t = nl for the dark rings, Eq. (vi) reads
rn2 = nlR.
or Dn2 = 4nlR
\ Dn µ n (ix)
Thus the diameter Dn of dark circular fringe(s) is proportional to the square root(s) of the natural numbers.
1.10.2.1 Determination of Wavelength of Light
We have seen that the diameter of nth order dark fringe in Newton’s rings method is
Dn2 = 4nl R (x)
From the above relation, the diameter of (n + p)th order dark fringe can be written as
D(2n + p ) = 4(n + p )l R (xi)
Subtracting Eq. (x) from equation (xi), we get
D(2n + p ) - Dn2 = 4 pl R
D(2n + p ) - Dn2
or l=
4 pR
Therefore, the measurement of diameters of the nth and (n + p)th dark fringes together with the radius of
curvature of the lens gives us the wavelength of sodium light with the help of above formula.
1.10.2.2 Determination of Radius of Curvature of Plano Convex Lens
This is clear from the theory of Newton’s rings that the measurement of diameters of nth order and (n + p)th
order dark fringes play an important role in the determination of wavelength of monochromatic light. For this
purpose, the following relation is used
D(2n + p ) - Dn2
l=
4 pR
Therefore, if we use the monochromatic source of light of known wavelength, it would be possible to
determine the radius of curvature of the plano convex lens with the help of following formula
D(2n + p ) - Dn2
R=
4 pl
1.10.2.3 Determination of Refractive Index of a Liquid

The liquid whose refractive index is to be determined is placed between the lens and the glass plate and then
we evaluate the diameters of the dark fringes.
The diameter of nth order dark fringe in air film is given by
Dn2 = 4nl R
Similarly, the diameter of nth order dark fringe in liquid film would be
4nl R
[ Dn2 ]liquid =
m
where m is the refractive index of the liquid and Dliquid < Dair
Therefore, the refractive index of the liquid can be calculated from the following formula once we are able to
determine the diameters of dark fringes.
S
[D2 ]
m = 2n air
[ Dn ]liquid
1.10.2.4 Newton’s Rings in Transmitted Light

Newton’s rings can be observed in reflected as well as in
transmitted light. Figure 1.22 shows that the rays QA and
HRB are the transmitted rays, which interfere. From the
figure it is also clear that the ray QA suffers no reflection
at a medium of higher index, so its phase does not change. H P
However, the ray HRB encounters two reflections at the
denser medium at Q and H. Since a phase change of p O Q
occurs at each reflection, the total phase change due to
both reflections would be 2p. Therefore, there will not
R
be any phase shift. In view of this, the path difference
Transmitted
between the two transmitted rays QA and HRB would be Rays
D = 2mt cos(r + q) (i)
For air (m = 1), normal incidence (r = 0) and smaller angle
q (cos q = 1), the path difference becomes B A
D = 2t (ii) Figure 1.22
The above equations shows that at t = 0, the path difference between the two transmitted rays D = 0. Therefore,
at the centre, the bright fringe will appear.
From Eq. (ii), the conditions for maxima and minima can respectively be obtained as below
2t = nl, n = 0, 1, 2, … (iii)
2t = (n + 1/2)l, n = 0, 1, 2,... (iv)
Because of the same reason as discussed earlier, the fringes in the transmitted light will also be circular. The
diameter of bright circular fringes can be obtained as
Dn = 2 nl R
Thus the diameter of the bright fringes is proportional to the square root of natural numbers. When we
calculate the diameter of dark circular fringes, it comes out to be
Dn = 2(n + 1)l R
This relation shows that the diameter of the dark fringes is proportional to the square root of odd natural
numbers.
Interference 27
From the above relation, it is clear that the fringes observed in the transmitted light are exactly complementary
to that of the reflected light. These fringes are much poorer in contrast as the transmitted rays emerge with
lower intensity in comparison with the reflected rays. The Newton’s rings obtained in the reflected as well as
in the transmitted light are shown in Fig. 1.23a and b, respectively.
(a) (b)
Figure 1.23
1.10.2.5 Newton’s Rings formed between Two Curved Surfaces

Let us consider two curved surfaces of radii of curvature R1 and R1 P
R2 in contact at the point O. A thin air film is enclosed between Air film
the two curved surfaces (Fig. 1.24). In this arrangement also, R2
Q
dark and bright rings are formed and can be seen with a
traveling microscope. O
rn
T
The thickness of the air film at P is Figure 1.24

PQ = PT – QT
If the radius of nth dark ring be rn, then from geometry,
rn2 r2
PT = and QT n
2 R1 2 R2
rn2 r2
\ PQ = - n
2 R1 2 R2
If we assume the thickness of the film as t, then
rn2 r2
t= - n
2 R1 2 R2
Now, this is clear from Fig. 1.24 that this type of film is similar to the wedge shaped film. Therefore, the path
difference between the wave reflected from the upper and lower surfaces of the film would be
l
2 mt cos (r + q ) +
2
For air (m = 1), normal incidence (r = 0) and the smaller angle q, the path difference takes the form
l
2t +
2
Therefore, in case of reflected light, for nth dark fringes
l Ê 1ˆ
2t + = Án + ˜ l
2 Ë 2¯
or 2t = nl
È r2 r2 ˘
or 2 Í n - n ˙ = nl
Î 2 R1 2 R2 ˚
È1 1˘
rn2 Í - ˙ = nl where n = 0, 1, 2, 3… (i)
Î R1 R2 ˚
Similarly, for nth bright fringe the path difference should satisfy the following condition
l
2t + = nl
2
Ê 1ˆ
or 2t = Á n - ˜ l
Ë 2¯
È1 1 ˘ (2n - 1)l
rn2 Í - ˙ = where n = 0, 1, 2, 3... (ii)
R
Î 1 R2˚ 2
Thus, bright and dark fringes are obtained according to Eqs. (i) and (ii). The diameter of the fringes can also
be calculated.
Now we invert the lower surface of the film. Under this R1
P
situation, the film would appear thicker than the previous case Air film
(Fig. 1.25). The film thickness PQ in this case would be
PQ = PT + QT T
O
rn2 r2 m=1
t= + n Q
2 R1 2 R2 R2 rn
For the reasons explained in wedge-shaped film, the following
condition should be satisfied in order to obtain nth order dark Figure 1.25
fringe of radius rn
2t = nl (for air)
È 1 1 ˘
or 2rn2 Í + ˙ = nl
Î 2 R1 2 R2 ˚
È1 1˘
rn2 Í + ˙ = nl where n = 0, 1, 2, 3... (iii)
R
Î 1 R2˚
Interference 29
For nth bright fringe

l
2t = (2n - 1)
2
È 1 1 ˘ l
or 2rn2 Í + ˙ = (2n - 1) 2
2
Î 1 R 2 R2˚
È1 1˘ l
rn2 Í + ˙ = (2n - 1) where n = 0, 1, 2, ....etc. (iv)
R
Î 1 R2˚ 2
A comparison of Eq. (i) with Eq. (iii) reveals that the diameter of dark fringes in the second case, where below
curved surface looks like convex when viewed from above, would be smaller than the one in first case. This
effect is similar to the situation as if we increase the width or thickness of the film. The same is the case for
bright fringes.
1.10.3 Michelson’s Interferometer
It consists of two highly polished mirrors M1 and M2 and M1 M2¢
two plane glass plates P and Q parallel to each other, as
shown in Fig. 1.26. The glass plate P is half-silvered on
its back surface and inclined at an angle of 45° to the L
beam of incident light. Another glass plate Q is such that
P and Q are of equal thickness and of the same material. S
45o
O
Two plane mirrors M1 and M2 are silvered on their front P Q
surfaces and mounted on two arms at right angle to each
M2
other. The position of the mirror M1 can be changed with
the help of a fine screw.
Light from a monochromatic source S, rendered parallel T
by a lens L, falls on the glass plate P. The semi-silvered Figure 1.26
plate P divides the incident light beam into two parts of
nearly equal intensities, namely reflected and transmitted beams. The reflected beam moves towards mirror
M1 and falls normally on it and hence it is reflected back to P and enters the telescope T. The transmitted
beam moves towards mirror M2 and falls normally on it after passing through the plate Q. Therefore, it is
reflected by the mirror M2 and follows the same path. At P it is reflected to enter the telescope T. Since the
beams entering the telescope have been derived from the same incident beam, these two rays are capable of
giving the phenomenon of interference; thereby producing interference fringes.
Function of Plate Q: The beam going towards the mirror M1 and reflected back, crosses the plate P twice,
while the other beam in the absence of Q would travel wholly in air. Therefore, to compensate the additional
path, the plate Q is used between the mirror M2 and plate P. The light beam going towards the mirror M2 and
reflected back towards P also passes twice through the compensation plate Q. Therefore, the optical paths of
the two rays in glass are the same.
Types of Fringes: The fringes in Michelson interferometer depend upon the inclination of M1 and M2.
Let M 2¢ be the image of M2 formed by the reflection at the half-silvered surface of the plate P so that
OM2 = OM 2¢ . The interference fringes may be regarded as formed by the light reflected from the surfaces of
M1 and M 2¢ . Thus, the arrangement is equivalent to an air-film enclosed between the reflecting surfaces M1
and M 2¢ .
It is obvious that the path difference between the two beams produced by the reflecting surfaces M1 and M 2¢
is equal to the twice of the thickness of the film M1 M 2¢ . This path difference can be varied by moving M1
backwards or forward parallel to itself. If we use monochromatic light, the pattern of bright and dark fringes
will be formed. Here the shape of the fringes will depend upon the inclination of M1 and M2.
If M1 and M2 are exactly at right angles to each other, the reflecting surfaces M1 and M 2¢ are parallel and hence
air film between M1 and M 2¢ is of constant thickness t so that we get circular fringes of equal inclination.
These fringes are called as Haidinger’s fringes that can be seen in the field view of a telescope. When the
distance between the mirrors M1 and M2 or between M1 and M 2¢ is decreased, the circular fringes shrink and
vanish at the centre. A ring disappears each time when the path 2t decreases by l.
Since the vertical ray first gets reflected from the inner surface of P (internal reflection), and then from the
front surface of the mirror M1 (external reflection) a phase change of p takes place. The horizontal ray first
gets reflected from the front surfaces of M2 (external reflection) and then from the inner surface of glass plate
P (external reflection), so there is no phase change. Therefore, the total path difference for normal incidence
would be
l
D = 2t cos q +
2
For bright fringes, the following condition should be satisfied
Ê 1ˆ [Q D = l]
2t cos q = Á n - ˜l (i)
Ë 2¯
For dark fringes, the condition reads
È Ê 1ˆ ˘
2t cos q = nl ÍQ D = ÁË n + 2 ˜¯ l ˙ (ii)
Î ˚
When t is further decreased, a limit is attained where M1 and M 2¢ coincide and the path difference between
the two rays becomes zero. Now the field of view is perfectly dark. When M1 is further moved, the fringes
appear again.
If M1 and M2 are not perfectly perpendicular, a wedge shaped film will be formed between M1 and M 2¢ then
we get almost straight line fringes of equal thickness in the field of view of telescope, as the radius of fringes
is very large.
All the above discussed films are shown in Fig. 1.27.
(a) (b) (c) (d)
Figure 1.27
Interference 31
1.10.3.1 Applications
Michelson’s interferometer uses the concept of interference that takes place with the help of two mirrors. The
distance between one mirror and the image of another plays an important role in the formation of fringes.
Michelson’s interferometer has diverse applications, some of which are listed below.
(i) Determination of Wavelength of Light
First of all the Michelson’s interferometer is set for circular fringes with central bright spot, which is possible
when both the mirrors are parallel (q = 0). If t be the thickness of air film enclosed between the two mirrors
(M1 and M 2¢ ) and n be the order of the spot obtained, then for normal incidence cos r =1, we have
l
2t + = nl
2
Ê 1ˆ
or 2t = Á n - ˜ l
Ë 2¯
l
If M1 is moved away from M 2¢ , then an additional path difference of l will be introduced and hence
2 l
(n +1)th bright spot appears at the centre of the field. Thus each time M1 moves through a distance , a new
2
bright fringe appears. Therefore, if M1 moves by a distance x (x1 to x2) and N new fringes appear at the centre
of the field, then we have
l
x = x2 - x1 = N
2
2( x2 - x1 ) 2 x 2x
or l= = l=
N N N
The difference (x2 – x1) is measured with the help of micrometer screw and N is actually counted. The
experiment is repeated for number of times and the mean value of l is obtained.
(ii) Determination of Difference in Wavelengths
Michelson’s interferometer is adjusted in order to obtain the circular fringes. Let the source be not
monochromatic and have two wavelengths l1 and l2 (l1 > l2) which are very close to each other (as Sodium
D lines). The two wavelengths form their separate fringe patterns but as l1 and l2 are very close to each other
and thickness of air film is small, the two patterns practically coincide with each other. As the mirror M1 is
moved slowly, the two patterns separate slowly and when the thickness of air film is such that the dark fringe
of l1 falls on bright fringe of l2, the result is maximum indistinctness. Now the mirror M1 is further moved,
say through a distance x, so that the next indistinct position is reached. In this position, if n fringes of l1
appear at the centre, then (n +1) fringes of l2 should appear at the centre of the field of view. Hence
l1 l
x=n and x = (n + 1) 2
2 2
2x
or n= (i)
l1
2x
and (n + 1) = (ii)
l2
On subtracting Eq. (i) from Eq. (ii), we get

2x 2x
(n + 1) - n = -
l2 l1
2 x(l1 - l2 )
or 1=
l1l2
2
l1l2 lav 2
or (l1 - l2 ) = = where l1 l2 = lav is the square of mean of l1 and l2.
2x 2x
Thus measuring the distance x moved by mirror M1 between the two consecutive positions of maximum
indistinctness, the difference between two wavelengths of the source can be determined, if lav is known.
(iii) Determination of Thickness and Refractive Index of a Thin Transparent Sheet
The Michelson’s interferometer is adjusted for producing straight white light fringes and cross-wire is set up
on the central bright fringe. Now insert thin transparent plate in the path of one of the interfering waves. On
the inclusion of a plate of thickness t and refractive index m, the path difference is increased by a factor of
2(m – 1)t. The fringes are therefore shifted. The mirror M1 is now moved till the central fringe is again brought
back to its initial position. The distance x traveled by the mirror M1 is measured by micrometer. Therefore
x
2 x = 2( m - 1)t or t = (iii)
( m - 1)
From Eq. (iii), we can write
x
m= +1 (iv)
t
Thus, by knowing the thickness of the transparent sheet and the distance x, we can calculate the refractive
index of the sheet with the help of a Michelson’s interferometer.
1.11 appLiCatiOns Of interferenCe in tHe fieLD Of engineering LO7

The phenomenon of interference arises in many situations and the scientists and engineers have taken
advantage of interference in designing and developing various instruments.
1.11.1 Testing of Optical Flatness of Surfaces

An example of the application of interference method is the testing of optical components for surface quality.
The most important example is that of optical flats. However, the methods used for fl at surfaces can be
adapted simply to test spherical surfaces.
1.11.1.1 Flatness Interferometers
With these interferometers we can compare the flatness of two surfaces by placing them in contact with slight
wedge of air between them. This gives a tilt and thus the fringes start originating like that of Newton’s ring
between the two surfaces. To get half wavelength contours of the space between the surfaces, they should be
viewed from infiity. Further, to avoid the risk of scratching, a desirable distance should be there between the
two surfaces. Most common examples of flatness interferometers are Fizeau and Twyman interferometers.
Interference 33
(i) Fizeau Interferometer

S
In this type of interferometer, the sources and viewing point are kept at
infinity (Fig. 1.28). This interferometer generates interference between
the surface of a test sample and a reference surface that is brought close
to the test sample. The interference images are recorded and analysed
by an imaging optic system. However, the contrast and the shape of the
interference signals depend on the reflectivity of the test samples.
(ii) Twyman-Green Interferometer
This is an important instrument used to measure defects in optical
components such as lenses, prisms, plane parallel windows, laser rods and
plane mirrors. Twyman-Green interferometer, shown in Fig. 1.29 resembles
Michelson interferometer in the beam splitter and mirror arrangement.
However, the difference lies in the way of their illumination. In the case
of Twyman-Green interferometer, we use a monochromatic point source Figure 1.28
which is located at the principal focus of a well-corrected lens whereas in Michelson interferometer an
extended source is used. If the mirrors M1 and M2 are perpendicular to each other and the beam-splitter BS
makes an angle of 45° with the normal of each mirror, then the interference is exactly analogous to thin film
interference at normal incidence. Therefore, completely constructive interference is obtained when d = ml/2,
where d is the path difference between the two arms adjusted by translating the mirror M1. The complete
destructive interference is obtained when d = (m + 1/2)l/2. With the help of rotation of mirror M2 we can see
fringes of equal thickness on the screen, as the angle of incidence is constant. This situation is analogous to
interference pattern observed with collimated light and a thin film with varying thickness. In order to test
the optical components, one of the mirrors is intentionally tilted to create fringes. Then the quality of the
component can be determined from the change in the fringe pattern when the component is placed in the
interferometer. Lens testing is specifically important for quantifying aberrations and measuring the focal
length.
Rotate M1
BS M2
Lens Translate
Screen
Figure 1.29
1.11.2 Nonreflecting or Antireflecting (AR) Coatings

Interference-based coatings were invented in November 1935 by Alexander Smakula, who was working for
the Carl Zeiss optics company. Antireflecting coatings are a type of optical coatings. These are applied to the
surface of lenses and other optical devices for reducing reflection. This way the efficiency of the system gets
improved since less light is lost. For example, in a telescope the reduction in reflections improves the contrast
of the image by elimination of stray light. In another applications a coating on eyeglass lenses makes the eyes
of the wearer more visible. The anti-reflecting coatings can be mainly divided into three groups.
1.11.2.1 Single-layer Interference Coatings
The simplest interference non-reflecting coating consists of a single quarter-wave layer of transparent material.
The refractive index of this material is taken to be equal to the square root of the substrate’s refractive index.
This theoretically gives zero reflectance at the center wavelength and decreased reflectance for wavelengths
in a broad band around the center. The use of an intermediate layer to form an antireflection coating can be
thought of as analogous to the technique of impedance matching of electrical signals. A similar method is
used in fibre optic research where an index matching oil is sometimes used to temporarily defeat total internal
reflection so that light may be coupled into or out of a fiber.
The antireflection coatings rely on an intermediate layer not only for its direct reduction of reflection
coefficient, but also use the interference effect of a thin layer. If the layer thickness is controlled precisely and
it is made exactly one quarter of the light’s wavelength (l/4), then it is called a quarter-wave coating (Fig.
1.30). In this case, the incident beam I, when reflected from the second interface, will travel exactly half its
own wavelength further than the beam reflected from the first surface. The two reflected beams R1 and R2
will destructively interfere as they are exactly out of phase and cancel each other if their intensities are equal.
Therefore, the reflection from the surface is suppressed and all the energy of the beam is propagated through
the transmitted beam T. In the calculation of the reflection from a stack of layers, the transfer-matrix method
can be used.
I λ
4
n0 n1 ns
l
R1 T
R2
Figure 1.30
Interference 35
1.11.2.2 Multilayer Coatings or Multicoating

Multiple coating layers can also be used for reflection reduction. It is possible if we design them such that
the reflections from the surfaces undergo maximum destructive interference. This can be done if we add
a second quarter-wave thick higher-index layer between the low-index layer (for example, silica) and the
substrate. Under this situation, the reflection from all three interfaces produces destructive interference
and antireflection. Optical coatings can also be made with near-zero reflectance at multiple wavelengths or
optimum performance at angles of incidence other than 0°.
1.11.2.3 Absorbing Antireflecting Coatings
Absorbing antireflecting coatings are an additional category of antireflection coatings. These coatings are
useful in situations where low reflectance is required and high transmission through a surface is unimportant
or undesirable. They can produce very low reflectance with few layers. They can often be produced more
cheaply or at greater scale than standard non-absorbing anti-reflecting coatings. In sputter deposition system
for such films, titanium nitride and niobium nitride are frequently used.
1.11.2.4 Practical Problems with AR Coatings
Real coatings do not reach perfect performance, though they are capable of reducing a surface’s reflection
coefficient to less than 0.1%. Practical details include correct calculation of the layer thickness. This is
because the wavelength of the light is reduced inside a medium and this thickness will be l0/4n1, where
l0 is the vacuum wavelength and n1 is the refractive index of the film. Finding suitable materials for use
on ordinary glass is also another difficulty, since few useful substances have the required refractive index
(n ª 1.23) which will make both reflected rays exactly equal in intensity. Since magnesium fluoride (MgF2)
is hard-wearing and can be easily applied to substrates using physical vapour deposition, it is often used for
this purpose even though its index is higher than desirable (n = 1.38).
1.12 sCientifiC appLiCatiOns Of interferenCe LO7

In interferometry, we use the principle of superposition to combine different waves in a way that will cause
the result of their combination to have some meaningful property, that is indicative of the original state of the
waves. The phenomenon of interference is employed under various situations for its scientific applications.
For a better understanding of the applications, we first need to know about the homodyne and heterodyne
detections.
1.13 HOmODYne anD HeterODYne DeteCtiOn LO7

In standard interferometry, the interference occurs between the two beams at the same wavelength (or carrier
frequency). The phase difference between the two beams results in a change in the intensity of the light on the
detector. Measuring the resulting intensity of the light after the mixing of these two light beams is known as
homodyne detection. In heterodyne detection, we modulate one of the two beams prior to detection, usually
by a frequency shift. A special case of heterodyne detection is optical heterodyne detection, which detects the
interference at the beat frequency.
1.13.1 Imaging Interferometry
In this interferometry, the pattern of radiation across a region can be represented as a function of position
i(x, y), i.e., an image and the pattern of incoming radiation i(x, y) can be transformed into the Fourier domain
f(u, v). A single detector measures information from a single point in (x, y) space. An interferometer measures
the difference in phase between two points in the (x, y) domain. This corresponds to a single point in the (u, v)
domain. An interferometer builds up a full picture by measuring multiple points in (u, v) space. The image
i(x, y) can then be restored by performing an inverse Fourier transform on the measured f(u, v) data.
1.13.2 Holographic Interferometry (HI)

Holographic interferometry (HI) is a technique that enables static and dynamic displacements of objects
with optically rough surfaces to be measured to optical interferometric precision, i.e., to fractions of a
wavelength of light. These measurements can be applied to stress, strain and vibration analysis, as well as
to nondestructive testing. It can also be used to detect optical path length variations in transparent media,
which enables, for example, fluid flow to be visualised and analysed. It can also be used to generate contours
representing the form of the surface. Holography interferometry is of two types.
(i) Live Holography Interferometry
Holography enables the light field scattered from an object to be recorded and replayed. If this recorded field
is superimposed on the “live field” scattered from the object, then the two fields will be identical. However,
if a small deformation is applied to the object, the relative phases of the two light fields will alter and it is
possible to observe interference. This technique is known as live holographic interferometry.
(ii) Frozen-Fringe Holography
In this holography, it is also possible to obtain fringes by making two recordings of the light field scattered
from the object on the same recording medium. The reconstructed light fields may then interfere to give
fringes, which map out the displacement of the surface.
1.13.3 Electronic Speckle Pattern Interferometry

Electronic Speckle Pattern Interferometry (ESPI), also known as TV Holography, is a technique that uses laser
light together with video detection, recording and processing to visualize static and dynamic displacements
of components with optically rough surfaces. The visualisation is in the form of fringes on the image where
each fringe normally represents a displacement of half a wavelength of the light used, i.e., quarter of a
micrometre or so.
1.13.4 Angle Resolved Low Coherence Interferometry

Angle resolved low coherence interferometry is an emerging biomedical imaging technology that uses the
properties of scattered light to measure the average size of cell structures, including the cell nuclei. The
technology shows promise as a clinical tool for in situ detection of dysplatic or precancerous tissue.
1.13.5 Optical Coherence Tomography

This is a medical imaging technique based on low-coherence interferometry, where subsurface light reflections
are resolved to give tomographic visualisation. Recent advances have struggled to combine the nanometre
phase retrieval with the ranging capability of low-coherence interferometry.
1.13.6 Geodetic Standard Baseline Measurements

A famous use of white light interferometry is the precise measurement of geodetic standard baselines. Here
the light path is split in two, and one leg is folded between a mirror pair 1 m apart. The other leg bounces once
off a mirror 6 m away. The fringes will be seen only if the second path is precisely 6 times the first. Starting
Interference 37
from a standard quartz gauge of 1 m length, it is possible to measure distances up to 864 m by repeated
multiplication. Baselines thus established are used to calibrate geodetic distance measurement equipments.
This leads to a metrologically traceable scale for geodetic networks measured by these instruments. More
modern geodetic applications of laser interferometry are in calibrating the divisions on levelling staffs
and in monitoring the free fall of a reflective prism within a ballistic or absolute gravimeter. This allows
determination of gravity, i.e., the acceleration of free fall, directly from the physical definition at a few parts
in a billion accuracy.
1.13.7 Interference Lithography
This is a technique for patterning regular arrays of fine features, without the use of complex optical systems
or photo marks. The basic principle of this is the same as interferometry. An interference pattern between two
or more coherent light waves is set up and recorded in a recording layer. This interference pattern consists of
a periodic series of fringes of representing intensity maxima and minima. The benefit of using interference
lithography is the quick generation of dense features over a wide area without loss of focus.
s UMMARY
We summarise the main outcome of the chapter as follows:
✦ We first discussed the phenomenon of interference and then explained it based on Young’s double slits
experiment.
✦ Concepts of wavefront and secondary wavelets were discussed based on Huygens’ principle. Then
secondary wavefront was introduced as the surface touching the secondary wavelets tangentially in the
forward direction at any given time.
✦ Phase difference and path difference between the two waves play a key role for obtaining constructive
or destructive interference. Therefore, phase and path differences were explained in detail together with
their relation.
✦ For obtaining sustained interference pattern, the two sources should be coherent. So the concept of
coherence, both temporal and spatial, was introduced and coherence time and coherence length were
talked about.
✦ A short description of a technique for producing coherent light from incoherent sources was given.
✦ Analytical treatment of the interference was discussed where conditions were obtained for the
constructive and destructive interferences.
✦ When a light wave gets reflected from a surface, a phase change may take place. Therefore, condition
of relative phase shift was explained.
✦ Superposition was extended for n number of waves and it was observed that the resultant amplitude
increases in proportion with n in length as n gets increased.
✦ Interfering waves can be produced either by division of wavefront or by division of amplitude.
Therefore, the details of interference were discussed based on these two methods.
✦ Fresnel’s biprism, which is used to create two virtual coherent sources, was discussed in detail for
obtaining interference pattern and the related conditions for dark and bright fringes.
✦ Application of biprism for the determination of wavelength of light, distance between two virtual
coherent sources and thickness of transparent sheet were discussed. The displacement of fringes by the
introduction of thin transparent sheet in the path of one light wave was also explored.
✦ Thin films are used for the division of amplitude of light waves which superimpose each other. The
interference pattern obtained by thin films of uniform and non-uniform thicknesses was investigated.
✦ When the air film is created between the curved surface of a plano-convex lens and the flat surface of
a mirror, the interference takes place between the reflected as well as the transmitted light. Here the
fringes are obtained in the form of rings known as Newton’s rings.
✦ Newton’s rings method was used for determination of the wavelength of light, radius of curvature of a
plano-convex lens and the refractive index of liquid.
✦ The theory was extended to Newton’s rings formed between two curved surfaces.
✦ Theory and practical applications of Michelson’s interferometer were discussed. Clarification of path
difference and the details of formation of fringes were given.
✦ Engineering applications of interference were included, particularly related to the testing of optical
flatness of surfaces and nonreflecting or antireflecting coatings.
✦ Finally the scientific applications of interference were discussed related to various interferometry,
tomography and lithography.
s olved e XAMPles
E xamplE 1 If light of wavelength 660 nm has wave trains 13.2 ¥ 10–6m long, what would be the coherence
time.
Solution Given l = 6.6 ¥ 10–7m, coherence length (DL) = 1.32 ¥ 10–5m and coherence time (Dt) = ?
Formula used is DL = c ◊ Dt
DL 1.32 ¥ 10-5
or Dt = = = 4.4 ¥ 10 –14 sec
c 3 ¥ 108
E xamplE 2 Coherence length of a light is 2.945 ¥ 10–2 m and its wavelength is 5896 Å. Calculate the
coherence time and number of oscillations corresponding to the coherence length.
Solution Given DL = 2.945 ¥ 10–2m and l = 5.896 ¥ 10–7 ◊ Dt = ?
Formula used is DL = cDt
DL 2.945 ¥ 10-2
or Dt = = = 9.816 ¥ 10 –11 sec
c 3 ¥ 108
and number of oscillations in a length L,
DL 2.945 ¥ 10-2
n= = = 4.99 ¥ 104
l 5.896 ¥ 10-7
E xamplE 3 A coherent beam has band width of 1200 Hz. Obtain the coherence length.
Solution Given Dn = 1200 Hz and C = 3 ¥ 108 m/s
Interference 39
C
By using the relation coherence length LC =
Dn
3.0 ¥ 108
LC = = 2.5 ¥ 105 m
1200
E xamplE 4 Calculate the line-width, coherence time and frequency stability for a line of Krypton having a
wavelength of 6.058 ¥ 10–7m and coherence length as 0.2 m.
Solution Given lav = 6.058 ¥ 10–7m, DL = 0.2 m and c = 3 ¥ 108m/sec.
In Michelson’s Interferometer we derived the following formula
2x 2x
1= -
l2 l1
where x is the distance between two mirrors. The above expression can be written as
2
l1l2 lav
2x = =
Dl Dl
where lav is the mean wavelength of l1 and l2. Here Dl is called line width.
In view of the fact that the fringes are not observed if the path difference exceeds the coherence length DL, we can assume
the beam to contain all wavelengths lying between l and (l + dl).
2
lav
Therefore, 2x = DL =
Dl
2
lav
or Dl =
DL
c
Q frequency, n =
l
È Dn ˘ c
\ = 2
ÎÍ Dl ˚˙ lav
c c
or Dn = 2 Dl =
lav DL
Here, Dv is called frequency spread of the line, which can be written in terms of Dt as follows.
1 1
Dv = or Dt =
Dt Dv
In addition, frequency stability is defined as the ratio of frequency spread and frequency of any spectral line, i.e.,
Dv
frequency stability =
v
2
lav (6.058 ¥ 10-7 )2
Line width Dl = = = 1.834 ¥ 10-12 m
DL 0.2
c 3 ¥ 108
Frequency spread Dv = = = 1.5 ¥ 109 Hz
DL 0.2
c 3 ¥ 108
Frequency v= =
l 6.058 ¥ 10-7
= 4.952 ¥ 1014 Hz
and Frequency stability
Dv 1.5 ¥ 109
= = = 3.0 ¥ 10 –6
v 4.952 ¥ 1014
E xamplE 5 The Doppler width for an orange line of Krypton is 550 ¥ 10–15 m. If the wavelength of light is
605.8 nm, calculate the coherent length.
Solution Given Doppler line width (Dl) = 5.5 ¥ 10–13m.
l = lav = 6.058 ¥ 10–7 m, DL = ?
2
lav (6.058 ¥ 10-7 )2
Formula used is DL = = = 0.6673 m
Dl 5.5 ¥ 10-13
E xamplE 6 A mercury vapour lamp emits a light of wavelength 5461 Å with a band width of 6 ¥ 108 Hz.
Calculate the ratio of its coherence length with the coherence length of a He–Ne laser operating at a wavelength
6328 Å with a band width of 106 Hz.
Solution For mercury vapour lamp, lav = 5.461 ¥ 10–7 m, = 6 ¥ 108 Hz
2
È Dv ˘ c lav Dv
Formula used is ÍÎ Dl ˙˚ = l 2 or Dl = c
av
(5.461 ¥ 10-7 ) 2 ¥ 6 ¥ 108

or Dl =
3 ¥ 108
= 5.964 ¥ 10-13 m
The coherence length is given by
2
lav (5.641 ¥ 10-7 ) 2
DL1 = = = 0.534 m
Dl 5.964 ¥ 10-13
For He – Ne laser,
Given lav = 6.328 ¥ 10–7 m, Dv = 106 Hz
2
lav Dv (6.328 ¥ 10-7 ) 2 ¥ 106
Dl = =
c 3 ¥ 108
–15
Dl = 1.335 ¥ 10 m
Coherence length (for laser)
2
lav (6.328 ¥ 10-7 )2
DL2 = = = 299.952 m
Dl 1.335 ¥ 10-15
DL1 0.534 1
\ = = = 1:562
DL2 299.952 562
E xamplE 7 Find the coherence length of a laser beam for which the band width is 3000 Hz.
Solution Given Dv = 3000 Hz.
1
Coherence length (DL) = cDt and coherence time (Dt) =
Dv
1
So Dt = = 3.333 ¥ 10-4 sec
3000
L = c ¥ t = 3 ¥ 108 ¥ 3.333 ¥ 10–4 = 1.0 ¥ 105m
E xamplE 8 Calculate the resultant line-width, band width and coherence length assuming that we chop a
continuous perfectly monochromatic beam of wavelength 6328 Å in 10–10 seconds using some sort of shutter.
Solution Given lav = 6.328 ¥ 10–7 m and t = 10–10 sec
Coherence length, DL = cDt = 3 ¥ 108 ¥ 10–10 = 3 ¥ 10–2 m
Interference 41
1 1
Bandwidth, Dg = = -10 = 1010 Hz
Dt 10
l 2 av (6.328 ¥ 10-7 )2 ¥ 1010
Line-width, Dl = Dg =
C 3 ¥ 108
= 1.335 ¥ 10–11 m
= 0.1335 Å
E xamplE 9 For a red cadmium line of wavelength 6438 Å and coherence length 38 cm deduce the order of
magnitude of (a) coherence time and (b) spectral width of the line.
Solution Given coherence length, DL = 0.38 m and lav = 6.43 8 ¥ 107 m
Coherence time Dt = ?
Spectral line width Dl = ?
DL
DL = c D t or Dt =
c
0.38
Dt = = 1.266 ¥ 10-9 sec
3 ¥ 108
2
lav l2 (6.438 ¥ 10-7 ) 2
or DL = or Dl = av =
Dl DL 0.38
Dl = 1.09 ¥ 10–12 m
E xamplE 10 The ratio of intensities of two waves that produce interference pattern is 16:1. Deduce the ratio
of maximum to minimum intensities in fringe system.
Solution Given I1 : I2 = 16 : 1
The intensity, I µ a2
\ a12 : a22 = 16 : 1 or a1 : a2 = 4 : 1
or a1 = 4a2
I max (a1 + a2 ) 2 (4a2 + a2 ) 2 25
= = =
I min (a1 - a2 ) 2 (4a2 - a2 ) 2 9
i.e., Imax : Imin = 25:9
E xamplE 11 Two waves of same frequency with amplitudes 1.0 and 2.0 units, interfere at a point, where the
phase difference is 60°. What is the resultant amplitude?
Solution Given a1 = 1.0 unit, a2 = 2.0 unit and f = 60°
the resultant amplitude
R = a12 + a22 + 2a1a2 cos f

= 1+ 4 + 2 = 7
= 2.65 units
E xamplE 12 Distance between two slits is 0.1 mm and the width of the fringes formed on the screen is 5 mm.
If the distance between the screen and the slit is one meter, what would be the wavelength of light used?
Solution Given b = 5.0 ¥ 10–3 m, 2d = 1.0 ¥ 10–4 m and D = 1.0 m

b 2d 5 ¥ 10-3 ¥ 1.0 ¥ 10-4
Formula used is l= = = 5000 Å
D 1.0
l = 5000 Å
E xamplE 13 A biprism of angle 1° and refractive index 1.5 is at a distance of 40 cm from the slit. Find the
fringe width at 60 cm from the biprism for sodium light of wavelength 5893 Å.
Solution Given a = .4 m, m = 1.5, l = 5.893 ¥ 10–7 m and D = 1.0 m
d
From the Fig. 1.31 d = or d = ad
a
or 2d = 2ad
The deviation produced in the incident light is given by
d = (m – 1) a
2d = 2a (m – 1)a
where a = angle of prism
p
1∞ = rad
180
lD 5.893 ¥ 10-7 ¥ 1.0 5.893 ¥ 10-7 ¥ 180
Fringe width b= = = = 0.0844 ¥ 10-3 m
2d 2 ¥ 0.4(1.5 - 1)p /180 2 ¥ 0.4 ¥ 0.5 ¥ 3.14
a d
d
2d
E xamplE 14 Interference fringes are produced by Fresnel’s bi-prism on the focal plane of a reading
microscope which is 1.0 m far from the slit. A lens interposed between the biprism and the microscope
gives two images of the slit in two positions. If the images of the slits are 4.05 mm apart in one position
and 2.90 mm apart in the other position and the wavelength of the sodium light is 5893 Å, find the distance
between the consecutive interference bands?
Solution Given l = 5.893 ¥ 10–7 m, D = 1.0 m, d1 = 4.05 ¥ 10–3 m and d2 = 2.90 ¥ 10–3 m
Formula used is 2d = d1d 2 = 4.05 ¥ 2.90 ¥ 10- 6
or 2d = 3.427 ¥ 10–3 m
Interference 43
l D 5.893 ¥ 10-7 ¥ 1.0

Now b= =
2d 3.427 ¥ 10-3
b = 0.172 mm
E xamplE 15 In a biprism experiment fringes were first obtained with the sodium light of wavelength 5890 Å
and fringe width was measured to be 0.342 mm. Sodium light was then replaced with white light and central
fringe was located. On introducing a thin glass sheet in half of the beam, the central fringe was shifted by
2.143 mm. Calculate the thickness of the glass sheet if the refractive index of glass is 1.542.
Solution Given l = 5890 ¥ 10–10 m, = 1.542, xn = 2.143 ¥ 10–3 m
b = 3.42 ¥ 10–4 m
Formula used is xn = nb
xn 2.143 ¥ 10-3
or n= = = 6.266
b 3.42 ¥ 10- 4
\ nª6
Q (m – 1)t = nl
nl 6 ¥ 5890 ¥ 10-10
or t= =
( m - 1) 0.542
t = 6.52 ¥ 10–6m.
E xamplE 16 Biprism is kept 10 cm away from the slit illuminated by monochromic light of l = 5896 Å. The
width of the fringes obtained on a screen placed at a distance of 90 cm from the biprism is 9.0 ¥ 10–4 m. What
is the distance between two coherent sources?
Solution Given a = 0.10 m, b = 0.90 m and D = a + b = 1.0 m
l = 5.896 ¥ 10–7 m, b = 9.0 ¥ 10–4 m
l D 5.896 ¥ 10-7 ¥ 1.0
Formula used is 2d = =
b 9 ¥ 10-4
or 2d = 6.55 ¥ 10–4m
E xamplE 17 The distance between the slit and biprism and between biprism and screen are 50 cm each,
Angle of biprism and refractive index are 179° and 1.5, respectively. Calculate the wavelength of light used
if the distance between two successive fringes is 0.0135 m.
Solution Given b = 0.0135 m, a = b = 0.5 m and D = a + b = 1
∞
180 - A Ê 1 ˆ p p
m = 1.5, A = 179∞, a = =Á ˜ ¥ = rad
2 Ë 2 ¯ 180 360
2a ( m - 1)a 2 ¥ 0.50 ¥ (1.5 - 1) p
Formula used is l= b= ¥ ¥ 0.0135
D 1.0 360
l = 5893 Å
E xamplE 18 The distance between the slit and biprism and between biprism and eyepiece are 45 cm each.
The obtuse angle of biprism is 178° and its refractive index is 1.5. If the fringe width is 15.6 ¥ 1 0–3 cm, find
the wavelength of light used.
lD b (2d )
Solution b= or l =
2d D
Given a = 45 cm, D = 90 cm, m = 1.5, a = 1° = p/180 rad
b = 15.6 ¥ 10–3 cm
2d can be calculated by the relation
2d = 2a (m – 1) a
= 2 ¥ 45 ¥ 0.5 ¥ (22/7) ¥ (1/180) = 0.786
b (2d ) 15.6 ¥ 10-3 ¥ 0.789
l= =
0 90
= 13676 Å
E xamplE 19 In a biprism experiment, the eye piece was placed at a distance of 120 cm from the source.
Calculate the wavelength of light, if the eye is required to move through a distance of 1.9 cm for 20 fringes
and distance between two slits is 0.06 cm.
Solution Given xn = 1.9 cm, n = 20, D = 120 cm and 2d = 0.06 cm.
xn 1.9 lD
Formula used is b= = = 0.095 cm and b =
n 20 2d
b 2d 0.095 ¥ 0.06
or l= =
D 120
l = 4750 Å
E xamplE 20 In a biprism experiment using light of wavelength 5890 Å, 40 fringes are observed in the field
of view. If this light is replaced by light of wavelength 4358 Å. Calculate how many fringes are observed in
the field of view.
Solution Given l1 = 5890 Å, N1 = 40 and l2 = 4358 Å, N2 = ?
l1D l D
x = N1b1 = N1 = N2 2
2d 2d
\ N1l1 = N2l2
40 ¥ 5890 ¥ 10–10 = N2 ¥ 4358 ¥ 10–10
N2 = 54
E xamplE 21 Light of wavelength 5893 Å is reflected at normal incidence from a soap film of refractive
index 1.42. What is the least thickness of the film that will appear (a) bright and (b) dark?
Solution Given l = 5.893 ¥ 10–7m, i = r = 0, = 1.42 and for smallest thickness n = 1
Condition for thin film to appear bright in reflected light is
2mt cos r = (2n – 1)l/2
(2n - 1)l / 2 (2 - 1) ¥ 5.893 ¥ 10-7
or t= =
2 m cos r 2 ¥ 1.42 ¥ 2 ¥ 1
= 1.038 ¥ 10 –4 mm
Similarly condition for thin film to appear to dark in reflected light is
2mt cos r = nl
Interference 45
nl 1 ¥ 5.893 ¥ 10-7
or t= =
2 m cos r 2 ¥ 1.42 ¥ 1
= 2.075 ¥ 10 –4 mm
E xamplE 22 A parallel beam of light strikes an oil film (m = 1.4), fl floating on a surface of water (m = 1.33).
When viewed at an angle of 30° from the normal 6th dark fringe is seen. Find the thickness of the film. (Given
wavelength of light 589 nm).
Solution Given l = 5.89 ¥ 10–7 m, moil = 1.4, i = 30° and n = 6
sin i sin 30∞ 0.5
moil = or 1.4 = or sin r =
sin r sin r 1.4
or sin r = 0.3571
or r = sin–1 (0.3571) = 20.92°
and cos r = cos 20.92 = 0.934
th
For n order dark fringe in reflected light, the condition is
nl
2 mt cos r = nl or t =
2 m cos r
6 ¥ 5.89 ¥ 10-7
t=
2 ¥ 1.4 ¥ 0.934
= 1.351 ¥ 10 –3 mm
E xamplE 23 Calculate the thickness of a soap film (m = 1.463) that will result in constructive interference
in the reflected light, if the film is illuminated normally with light whose wavelength in free space is 6000 Å.
Solution Given l = 6.0 ¥ 10–7 m, m = 1.463, for normal incidence i = r = 0° and for smallest thickness n = 1.
For constructive interference 2 mt cos r = (2n – 1) l/2
(2n - 1)l (2 - 1) ¥ 6.0 ¥ 10-7
t= =
2 ¥ 2 ¥ 1.463 ¥ 1 4 ¥ 1.463
= 1.025 ¥ 10 –4 mm
E xamplE 24 A parallel beam of sodium light (l = 5890 Å) strikes a film of oil floating on water. When
viewed at an angle of 30° from the normal, qth dark band is seen. Determine the thickness of the film.
(Refractive index of oil = 1.46).
Solution Given l = 5.89 ¥ 10–7 m, i = 30°, m = 1.46 and n = 8
Condition for obtaining dark band is 2mt cos r = nl (i)
nl
or t=
2 mt cos r
sin i
As we know, m= (ii)
sin r
sin i
or r= (iii)
m
sin 30∞ 1
or sin r = =
1.46 2.92
or cos r = 1 - sin 2 r
2
By using Eq. (ii), we get = 1 - ÊÁ
1 ˆ
Ë 2.92 ˜¯
= 0.94
8 ¥ 5.89 ¥ 10-7
t=
2 ¥ 1.46 ¥ 0.94
= 1.72 ¥ 10 –3 mm
E xamplE 25 White light is reflected from an oil film of thickness 0.01 mm and refractive index 1.4 at an
angle of 45° to the vertical. If the reflected light falls on the slit of a spectrometer, calculate the number of
dark bands seen between wavelengths 4000 and 5000 Å.
Solution Given t = 1.0 ¥ 10–5 m, m = 1.4, i = 45°, l1 = 4.0 ¥ 10–7 m and l2 = 5.0 ¥ 10–7 m
Condition of dark bands in reflected light is
2mt cos r = nl (i)
sin i sin i
m= or sin r =
sin r m
sin 45∞ 1/ 2 1
or sin r = = =
1.4 1.4 1.4 2
1
cos r = 1 - sin 2 r = 1 -
2 ¥ 1.96
= 0.86
For wavelength l1, i.e, 4.0 ¥ 10–7 m

2mt cos r = n1l1
2mt cos r 2 ¥ 1.4 ¥ 1.0 ¥ 10-5 ¥ 0.86
n1 = = = 60.2
l1 4.0 ¥ 10-7
n1 ª 60
For wavelength l2, i.e, 5.0 ¥ 10–7 m
2 mt cos r 2 ¥ 1.4 ¥ 1.0 ¥ 10-5 ¥ 0.86
n2 = =
l2 5.0 ¥ 10-7
= 48.16 or n2 ª 48
n1 – n2 = 60 – 48 = 12
i.e, 12 dark bands are seen between wavelengths 4000 and 5000 Å.
E xamplE 26 A parallel beam of light of wavelength 5890 Å is incident on a glass plate having refractive
index m = 1.5 such that the angle of refraction in the plate is 60°. Calculate the smallest thickness of glass
plate which will appear dark by reflected light.
Solution Given l = 5.89 ¥ 10–7 m, = 1.5 and r = 60°
Condition for the film to appear dark in reflected light is
2mt cos r = nl
Interference 47
For minimum thickness n = 1

l 5.89 ¥ 10-7
t= =
2 m cos r 2 ¥ 1.5 ¥ 0.5
= 0.3927 ¥ 10 –6 m
E xamplE 27 A soap film of refractive index 1.333 is illuminated by white light incident at an angle of
45°. The light refracted by it is examined by a spectroscope and two consecutive bright bands are focused
corresponding to the wavelength 6.1 ¥ 10–5 cm and 6.0 ¥ 10–5 cm. Find the thickness of the film.
Solution Given m = 1.333, i = 45°, l1 = 6.1 ¥ 10–7 m and l2 = 6.0 ¥ 10–7 m
sin i sin 45∞ 1/ 2 0.707
m= or sin r = = =
sin r 1.333 1.333 1.333
sin r = 0.53
cos r = 1 - sin 2 r = 1 - (0.53) 2 = 0.848

Condition of bright film to observe in transmitted case is
2mt cos r = nl1 = (n + 1)l2
or n ¥ 6.1 ¥ 10–7 = (n + 1) ¥ 6.0 ¥ 107
n = 60
nl1 60 ¥ 6.1 ¥ 10-7
and t= =
2 m cos r 2 ¥ 1.333 ¥ 0.848
= 1.62 ¥ 10 –5 mm
E xamplE 28 A soap film suspended in air has thickness 5 ¥ 10–5 cm viewed at an angle 35° to the normal.
Find the wavelength of light in visible spectrum, which will be absent for a reflected light. The m for the soap
film is 1.33 and visible spectrum is in the range of 4000 to 7800 Å
Solution Given t = 5.0 ¥ 10–7 m, i = 35° and m = 1.33
By using the relation
sin i
2mt cos r = nl and m =
sin r
sin i
sin r = and cos r = 1 - sin 2 r
sin m
2 2
Ê sin i ˆ Ê sin 35 ˆ
cos r = 1- Á = 1- Á
Ë m ˜¯ Ë 1.33 ˜¯
= 0.902
for first order (n = 1)
1 ◊ l1 = 2m + cos r
l1 = 2 ¥ 1.33 ¥ 5.0 ¥ 10–7 ¥ 0.902
= 1.19 ¥ 10–6
= 1200 Å
In second order (n = 2)
2 ¥ l2 = 2mt + cos r
or l2 = m + cos r = 1.33 ¥ 5 ¥ 10–7 ¥ 0.902

= 6.0 ¥ 10–7 m (approx)
3 ¥ l3 = 2mt + cos r
2 ¥ 1.33 ¥ 5.0 ¥ 10-7 ¥ 0.902
or l3 =
3
= 4000 Å (approx)
4 ¥ l4 = 2mt + cos r
2 ¥ 1.33 ¥ 5.0 ¥ 10-7 ¥ 0.902
or l4 =
4
= 3000 Å (approx)
Hence l2 and l3 wavelength of light in a visible spectrum will be absent.
E xamplE 29 Calculate the thickness of a soap film (m = 1.463) that will result in constructive interference
in the reflected light, if the film is illuminated normally with light whose wavelength in free space is 6000 Å.
Solution Given l = 6.0 ¥ 10–7 m, in this case n = 1, and for normal incidence i = 0 and r = 0
Condition for constructive interference 2mt cos r = (2n – 1) l/2
(2n - 1)l 1 ¥ 6.0 ¥ 10-7
t= =
4 m cos r 4 ¥ 1.463 ¥ 1
= 1.025 ¥ 10 –7 m.
E xamplE 30 A thin film is illuminated by white light at an angle of incidence (i = sin–1(4/5). In reflected
light, two dark consecutive overlapping fringes are observed corresponding to wavelengths 6.1 ¥ 10–7 m and
6.0 ¥ 1 0–7 m. The refractive index of the film is 4/3. Calculate the thickness of the film.
Solution Given l1 = 6.1 ¥ 10–7 m, l2 = 6.0 ¥ 10–7 m and m = 4/3
sin i sin i 4/5 3
m= or sin r = = =
sin r m 4/3 5
9 4
and cos r = 1 - sin 2 r = 1 - = = 0.8
25 5
Condition for dark fringes is
2mt cos r = nl1 = (n + 1)l2
or nl1 = (n + 1)l2
n ¥ 6.1 ¥ 10–7 = (n + 1) 6.0 ¥ 10–7
n(6.1 – 6.0) ¥ 10–7 = 6.0 ¥ 10–7
or n = 60
and 2mt cos r = nl1
60 ¥ 6.1 ¥ 10-7
or t=
4
2 ¥ ¥ 0.8
3
t = 1.716 ¥ 10 –2 mm
Interference 49
E xamplE 31 Two plane glass surfaces in contact along one edge are separated at the opposite edge by a thin
wire. If 20 interference fringes are observed between these edges in sodium light at normal incidence. What
is the thickness of wire?
Solution Given lav = 5893 Å, n = 20, i = r = 0 and m = 1
l 5.893 ¥ 10-7 m
w= =
2 mq 2 ¥1¥q
5.893 ¥ 10-7
or wq =
2 B
From Fig. 1.32
t
q=
20w t
or t = 20 wq
q
5.893 ¥ 10-7 O
= 20 ¥ 20 w A
2
Figure 1.32
t = 5.893 ¥ 103 mm.
E xamplE 32 A wedge air film is enclosed between two glass plates touching at one edge and separated by a
wire of 0.06 ¥ 10–3 m diameter at a distance of 0.1 5 m from the edge. Calculate the fringe width. The light
of wavelength 6.0 ¥ 10–7 m from the broad source is allowed to fall normally on the film.
Solution Given l = 6.0 ¥ 10–7 m and m = 1 B
l
w= (i)
2 mq
From Fig. 1.33 6.0 × 10−5
6.0 ¥ 10-5 (ii)
q=
0.15
6.0 ¥ 10-7 ¥ 0.15 q
\ w= 0.15 A
2 ¥ 1 ¥ 6.0 ¥ 10-5
Figure 1.33
w = 0.75 mm.
E xamplE 33 A wedge shaped film is illuminated by light of wavelength 4650 Å. The angle of wedge is 40≤.
Calculate the fringe separation between two consecutive fringes.
Solution Given l = 4.56 ¥ 10–7 m and m = 1
40 p
q = 40¢¢ = ¥ rad
3600 180∞
= 1.9 ¥ 10-4 rad
l 4.65 ¥ 10-7
\ w= =
2 ma 2 ¥ 1 ¥ 1.9 ¥ 10-4
w = 1.2 mm
E xamplE 34 Two glass plates enclose a wedge-shaped air film touching at one edge are separated by a wire
of 0.03 mm diameter at distance 15 cm from the edge. Monochromatic light (l = 6000 Å) from a broad source
falls normally on the film. Calculate the fringe-width.
Solution Given l = 6.0 ¥ 10–7 m and m = 1

È Arc (AB ) ˘
ÍÎAngle q = radius ˙˚
0.03 ¥ 10-3
q=
0.15
= 2.0 ¥ 10-4 rad
l
w=
2 mq
l
=
2q
6.0 ¥ 10-7
w= = 1.5 ¥ 10-3 m = 1.5 mm
2 ¥ 2 ¥ 10-4
E xamplE 35 A glass wedge having angle 0.01 radian is illuminated normally by light of wavelength 5890 Å.
At what distance from the edge of the wedge, will the 12th dark fringe be observed by reflected light?
Solution Given l = 5.89 ¥ 10–7 m, n = 12, q = 0.01 rad and m = 1
Condition for obtaining dark fringe is
2mt cos (r + q) = nl (i)
For normal incidence i = r = 0 and when q is very small
cos q ª 1
Eq. (i) reads 2t = nl (ii)
t
Now the angle q can be written as q =
x
where t is the thickness and x is the distance from the edge (Fig. 1.34) then
we have t=q◊x (iii)
t
By using Eqs. (ii) and (iii), we get
2q ◊ x = nl q
x
nl 12 ¥ 5.89 ¥ 10-7
or x= = = 3.5 ¥ 10–4 m
2q 2 ¥ 0.01 Figure 1.34
x = 0.35 mm.
E xamplE 36 A glass wedge of angle 0.01 radian of illuminated by monochromatic light of wavelength 6000
Å falling normally on it. At what distance from the edge of the wedge will the 10th fringe be observed by
reflected light?
Solution Given a = 0.01 radian, l = 10 ¥ l = 6.0 ¥ 10–7 m the condition for dark fringe
= 2t = nl (i)
t
The angle of wedge a =
x
or t = ax (ii)
Put the value of t from Eq. (ii) in Eq (i), we get
2ax = nl
nl 10 ¥ 6.0 ¥ 10-7
x= = = 3 ¥ 10–4 m
2a 2 ¥ 0.01
Interference 51
E xamplE 37 Interference fringes are produced when monochromatic light is incident normally on a thin
wedge-shaped film of refractive index 1.5. If the distance between two consecutive fringes is 0.02 mm. Find
the angle of the film, the wavelength of light being 5.5 ¥ 10–5 cm.
Solution Given m = 1.5, w = 0.02 ¥ 10–3 m and l = 5.5 ¥ 107 m.
l l 5.5 ¥ 10-7
w= or q = =
2 mq 2 m w 2 ¥ 1.5 ¥ 0.02 ¥ 10-3
= 0.009166 rad = 0.525°
E xamplE 38 In Newton’s rings experiment, the diameter of the 15th ring was found to be 0.59 cm and that
th
of the 5 ring was 0.336 cm. If the radius of the plano convex lens is 100 cm, compute the wavelength of
light used.
Solution Given D15 = 5.9 ¥ 10–3 m, D5 = 3.36 ¥ 10–3 m, p = 10 and r = 1.0 m.
D 2 (n + p ) - Dn2 [(5.9) 2 - (3.36) 2 ] ¥ 10-6
Formula used is l= =
4 pR 4 ¥ 10 ¥ 1.0
l= 5880 Å
E xamplE 39 In a Newton’s rings experiment the radius of 10th and 20th rings are 0.2 and 0.3 cm, respectively,
and the focal length of the plano-convex lens is 90 cm. Calculate the wavelength of light used in nanometers.
Solution Given f = 0.9 m, m = 1.5, D10 = 0.2 cm and D20 = 0.3 cm. p = 10
1 Ê 1 1ˆ È1 1˘
Formula used is = ( m - 1) Á - ˜ = (1.5 - 1) Í - ˙ , R = •
f Ë 1
R R2¯ R
Î 2 • ˚
1 È1˘
= 0.5 Í ˙ or R1 = R = 0.45 m
0.9 Î R1 ˚
D(2n + p ) - Dn2 2
D15 - D52 [(0.3) 2 - (0.2)2 ] ¥ 10-4
and l= = =
4 pR 4 ¥ 10 ¥ 0.45 4 ¥ 10 ¥ 0.45
l = 277.8 nm
E xamplE 40 In a Newton’s rings arrangement a thin convex lens of focal length 1.0 m. (m = 1.5) remains
in contact with an optical flat and light of wavelength 5896 ¥ 10–10 m is used. Newton’s rings are observed
normally by reflected light. What is the diameter of 7th bright ring?
Solution Given m = 1.5, f = 1.0 m and l = 5.896 ¥ 10–7 m
R1 = R and R2 = R
1 Ê 1 1ˆ
Formula used is = ( m - 1) Á - ˜
f Ë R1 R2 ¯
1 È1 1˘
or = (1.5 - 1) Í + ˙
1.0 ÎR R˚
2 1
or = or R = 1.0 m
R 0.5
now Dn2 = 4nl R
for n = 7
D7 = 4 ¥ 7 ¥ 5.896 ¥ 10-7 ¥ 1.0
D7 = 4.063 ¥ 103 m
E xamplE 41 Light source emitting the light of wavelengths l1 = 6.0 ¥ 10–7 m and l2 = 4.8 ¥ 10–7 m is used
to obtain Newton’s rings in reflected light. It is found that the nth dark ring of l1 coincides with (n +1)th dark
ring of l2. If the radius of curvature of the curved surface of the lens is 0.96 m. Calculate the diameter of
(n + 1)th dark ring of l2.
Solution Given l1 = 6.0 ¥ 10–7 m, l2 = 4.8 ¥ 10–7 m and R = 0.96 n
The diameter of nth order dark ring of l1 is
Dn2 (l1 ) = 4nl1R
Similarly, the diameter of (n + 1)th order dark ring of l2
D(2n +1) (l2 ) = 4(n + 1)l2 R
Since, Dn2 (l1 ) = D(2n +1) (l2 )

4nl1 = 4(n + 1)l2
n + 1 l1
or =
n l2
1 l1 1 l - l2
or 1+ = or = 1
n l2 n l2
l2 4.8 ¥ 10-7
or n= = = 4.0
l1 - l2 (6.0 - 4.8) ¥ 10-7
or n=4
Hence, D(2n +1) (l2) = 4(n + 1)l2R = 4 ¥ 5 ¥ 4.8 ¥ 10–7 ¥ 0.96
D(n + 1) = 3.0358 ¥ 10–3
or D(n+1) = 3.04 ¥ 10–3 m.
E xamplE 42 In Newton’s ring arrangement a source is emitting two wavelengths l1 = 6.0 ¥ 10–7 m and
l2 = 5.9 ¥ 10 m. It is found that nth dark ring due to one wavelength coincides with (n + 1)th dark ring due
–7
to the other. Find the diameter of the nth dark ring if radius of curvature of the lens is 0.9 m.
Solution Given l1 = 6.0 ¥ 10–7 m, l2 = 5.9 ¥ 10–7 m and R = 0.9 m.
The diameter of the nth order dark ring of l1 is
Dn2 (l1 ) = 4nl1R
The diameter of the (n + 1)th order dark ring of l2 is
D(2n +1) (l2 ) = 4( n +1) l2 R
Since two rings coincide
4nl1R = 4(n + 1)l2R
n + 1 l1 l2
= or n =
n l2 l1 - l2
5.9 ¥ 10-7
n= = 59
(6.0 - 5.9) ¥ 10-7
Now Dn = 4nl1 R = 4 ¥ 59 ¥ 6.0 ¥ 10-7 ¥ 0.9

= 0.01128 m
Dn = 0.0113 m
Interference 53
E xamplE 43 Newton’s rings are formed using light of wavelength 5896 Å in reflected light with a liquid
placed between plane and curved surfaces. The diameter of 7th bright fringe is 0.4 cm and the radius of
curvature is 1.0 m. Evaluate the refractive index of liquid.
Solution Given D7 = 4.0 ¥ 10–3 m, l = 5.896 ¥ 10–7 m, R = 1.0 m and n = 7.
2(2n - 1)l R 2(2n - 1)l R
Dn¢ 2 = or m =
m Dn¢ 2
2 ¥ 13 ¥ 5.896 ¥ 10-7 ¥ 1.0
m=
(4 ¥ 10-3 ) 2
m = 0.96
E xamplE 44 If the diameter of nth dark ring in an arrangement giving Newton’s ring changes from 0.3 cm
and 0.25 cm as liquid is introduced between the lens and the plate, calculate the value of the refractive index
of the liquid and also calculate the velocity of light in the liquid. Velocity of light in vacuum is 3 ¥ 108 m/sec.
Solution Given Dn = 3.0 ¥ 10–3 m, Dn = 2.5 ¥ 10–3 m
Formula used is Dn2 = 4nl R
4nl R
Dn2 =
m
2
D 2 È 3.0 ¥ 10-3 ˘
m = n2 = Í ˙ = 1.44
Dn Î 2.5 ¥ 10-3 ˚
m = 1.44
c 3 ¥ 108
or Vliq = =
m 1.44
Vliq = 2.08 ¥ 108 m/sec
E xamplE 45 The Newton’s rings are seen in reflected light of wavelength 5896Å. The radius of curvature of
plano-convex lens is 1.0 meter. An air film is replaced by a liquid whose refractive index is to be calculated
under the conditions if 16th ring is dark and its diameter is 5.1 mm.
Solution Given D16 = 5.1 ¥ 10–3 m, l = 5.896 ¥ 10–7 m and R = 1.0 m
4nl R 4nl R
Formula used is Dn2 = or
m Dn2
4 ¥ 16 ¥ 5.896 ¥ 10-7 ¥ 1.0
m=
(5.1 ¥ 10-3 )2
m = 1.45
E xamplE 46 The Newton’s rings are observed in reflected light of wavelength 6300 Å. A thin layer of liquid
of refractive index 1.63 is formed between curved surface of plano-convex lens (m = 1.69) and plane glass
plate (m = 1.03) and the radius of curvature of the curvex lens is 0.9 m. Find the radius of smallest dark ring.
Solution Given l = 6.3 ¥ 10–7 m, m = 1.63 and R = 0.9 m
nl R
Formula used is rn¢ 2 = (n = 1 for a smallest dark ring)
m
1 ¥ 6.3 ¥ 10-7 ¥ 0.9
rn¢ 2 = = 34.7853 ¥ 10-8 m
1.63
r12 = 5.9 ¥ 10–4 m = 0.59 mm
E xamplE 47 Newton’s rings are observed with two different media between the glass surfaces. The nth rings
have diameters as 10 : 7. Find the ratio of the refractive indices of the two media.
Solution Given Dn¢ : Dn≤ = 10 : 7
4nl R
Q Dn2 = (i)
m
For the first medium (m1)
4nl R
Dn¢ 2 = (ii)
m1
For the first medium (m1)
4nl R (iii)
Dn¢¢ 2 =
m2
2
m1 Dn¢¢ 2 È 7 ˘ 49
= = =
m 2 Dn¢ 2 ÎÍ10 ˚˙ 100
or m1 : m2 = 49 : 100
E xamplE 48 A combination of convex lens plane glass plate is illuminated by monochromatic light. The
diameter of the 10th dark ring is measured in reflected light and is found to be 0.48 cm. Find the wavelength
of light used. The radius of curvature of the lower face of the lens is 90 cm.
Solution Given R = 0.9 m, D10 = 4.8 ¥ 103 m and n = 10
Formula used is Dn2 = 4nl R
Dn2
or l=
4nR
(4.8 ¥ 10-3 ) 2
or l=
4 ¥ 10 ¥ 0.9
l = 6400 Å
E xamplE 49 In Newton’s rings experiment the diameter of 5th dark ring is reduced to half of its value after
placing a liquid between plane glass plate and convex surface. Calculate the refractive index of liquid.
D5
Solution Given D¢s = ◊m=?
2
2 2
Formula used is Dn = 4nlR = D5 = 4 ¥ 5 ¥ lR or D5 = 20l R
4 ¥ 5 ¥ lR 20l R
\ D5¢ 2 = or D5¢ =
m m
D5 20l R 20l R
D5¢ = or =
2 m 2
1 1
or = or m = 4
m 2
or m=4
E xamplE 50 Newton’s rings by reflection are formed between two bi-convex lenses having equal radii of
curvatures as 100 cm each. Calculate the distance between the 5th and 15th dark rings, using monochromatic
light of wavelength 5400 Å.
Interference 55
Solution Given l = 5.4 ¥ 10–7 m, R1 = R2 = 1.0 m

nl 5 ¥ 5.4 ¥ 10-7 ¥ 100
Formula used is rn2 = or r5 = = 11.62 ¥ 10- 4 m
Ï1 1 ¸ 1 1
Ì + ˝ +
R R 2 2
Ó 1 2˛
or r5 = 0.1162 cm.
15 ¥ 5.4 ¥ 10- 7
Similarly r15 = = 2.012 ¥ 10–3 m = 0.2012 cm
2
\ Distance between 5th and 15th rings = r15 – r5 = 0.085 cm
E xamplE 51 A Michelson interferometer is set for the white straight fringes. When a mica sheet of thickness
0.005 cm is put in front of the fixed mirror, then in order to bring back the coloured fringes to their original
position, the movable mirror is moved by 0.0025 cm. Calculate the refractive index of mica.
Solution Given x = 2.5 ¥ 10–5 m and t = 5.0 ¥ 10–5 m
Formula used is 2x = 2(m – 1)t
where t is the thickness of mica sheet and is the refractive index
x
or m = +1
t
2.5 ¥ 10-5 m
m= + 1 = 1.5
5.0 ¥ 10-5
or m = 1.5
E xamplE 52 If a movable mirror of Michelson’s interferometer is moved through a distance 0.06 mm, 200
fringes crossed the field of view. Find the wavelength of light.
Solution Given x = 6.0 ¥ 10–5 m and N = 200 fringes
Formula used is x = Nl/2
where x is the separation of movable mirror from the fixed mirror, then 6.0 ¥ 10–5 = 200 ¥ l/2 or = 6000 Å
E xamplE 53 In Michelson’s interferometer a thin plate is introduced in the path of one of the beams and it
is found that 50 band crosses the line of observation. If the wavelength of light used is 5896 Å and m = 1.4,
determine the thickness of the plate.
Solution Given n = 50, l = 5896 ¥ 10–10 m and m = 1.4
Formula used is 2(m – 1) t = nl
nl 50 ¥ 5896 ¥ 10-10
or = =
2( m - 1) 2(1.4 - 1)
t = 3.68 ¥ 10–5 m
E xamplE 54 Calculate the distance between successive positions of the movable mirror of Michelson’s
interferometer giving best fringes in case of a sodium source having wavelengths 5896 Å and 5890 Å. What
will be the change in path difference between two successive reappearances of the interference pattern?
Solution Given l1 = 5.896 m and 10–7l2 = 5.89 ¥ 10–7 m
l1l2
Formula used is Dl = (l1 - l2 ) =
2x
l1l2 5.896 ¥ 5.89 ¥ 10-14

or x= =
2(l1 - l2 ) 2 ¥ 6 ¥ 10-10
= 0.289 mm
The path difference will be equal to
= 2x = 2 ¥ 0.289 mm
= 0.5788 mm.
E xamplE 55 In Michelson’s interferometer 100 fringes cross the field of view when the movable mirror is
displaced through 0.02948 mm. Calculate the wavelength of monochromatic light used.
Solution Given x = 0.02948 ¥ 10–3 m and n = 100
2 x 2 ¥ 2.948 ¥ 10-5 m
Formula used is 2 x = nl or l = =
n 100
or l = 5.896 ¥ 10–7 m = 5896 Å
E xamplE 56 The wavelength of two components of D-lines of sodium are 5890 Å and 5896 Å. By how
much distance one of the mirror of Michelson’s interferometer be moved so as to obtain consecutive position
of maximum distinctness.
Solution Given l1 = 5.896 ¥ 10–7 m and l2 = 5.89 ¥ 10–7 m
l1l2
Formula used Dl = l1 - l2 =
2x
where x is distance through which the movable mirror is moved from one position of maxima to the next, then we have
l1l2 5.896 ¥ 5.89 ¥ 10-14
x= =
2(l1 - l2 ) 2 ¥ 6 ¥ 10-10
x = 0.289 mm
E xamplE 57 In an experiment with Michelson’s interferometer, the distance traveled by the mirror for two
successive position of maximum distinctness was 0.2945 mm. If the mean wavelength for the two component
of sodium D–line is 5893 Å, calculate the difference between the two wavelengths.
Solution Given x = 0.2945 ¥ 10–3 m and lav = 5.893 ¥ 10–7 m
2
lav (5.893 ¥ 10-7 ) 2
Formula used is Dl = l1 - l2 = =
2 x 2 ¥ 0.2945 ¥ 10-3
Dl = 5.896 Å.
E xamplE 58 In an experiment for determining the refractive index of a gas using Michelson’s interferometer
a shift of 140 fringes is observed when all the gas is removed from the tube. If the wavelength of light used
is 5460 Å and the length of the tube is 20 cm, calculate the refractive index of the gas.
Solution Given, l = 5.46 ¥ 10–7 m, t = 0.2 m and n =140
Formula used 2(m – 1) t = nl
nl 140 ¥ 5.46 ¥ 10-7
or m= +1= +1
2t 2 ¥ 0.2
= 0.00019 + 1
m = 1.00019
Interference 57
o BJeCtive t YPe Q Uestions
Q.1 Light waves are

(a) longitudinal waves (b) transverse waves
(c) both (a) and (b) (d) none of them
Q.2 Which of the following does not support the wave nature of light?
(a) interference (b) polarisation (c) compton effect (d) diffraction
Q.3 Colours in thin films are because of
(a) dispersion (b) diffraction (c) interference (d) none of them
Q.4 A phase difference p between two interfering beams is equivalent to path deference
(a) 2l (b) l (c) l/2 (d) none of them
Q.5 The central fringe in Fresnel’s biprism experiment in the reflected light is
(a) dark (b) non-uniform (c) bright (d) none of them
Q.6 In reflected light the central fringe of Newton’s ring is
(a) dark (b) bright (c) non-uniform (d) none of them
Q.7 In Newton’s rings experiment, the diameter of bright rings is proportional to
(a) odd natural numbers (b) natural numbers
(c) even natural numbers (d) square root of natural numbers
Q.8 Extended source is needed in
(a) Young’s double slit experiment (b) Biprism experiment
(c) Newton’s rings experiment (d) none of them
Q.9 Second glass plate in Michelson’s interferometer is known as
(a) extra glass plate (b) compensating glass plate
(c) simple glass plate (d) none of them
Q.10 Intensity of light is proportional to
(a) double of its amplitude (b) square root of amplitude
(c) square of amplitude (d) none of these
Q.11 Which of the following device cannot produce coherent beam of light?
(a) Fresnel’s biprism (b) Young’s double slit
(c) Half-shade polarimeter (d) Lloyd’s mirror
Q.12 Which of the following method(s) is/are based on division of amplitude?
(a) Newton’s rings (b) Michelson interferometer
(c) (a) and (b) (d) None of these
Q.13 Which of the following method(s) is(are) based on division of wave front?
(a) Young’s double slit (b) Newton’s rings
(c) Fresnel’s biprism (d) (a) and (c)
Q.14 What is the actual shape of interference fringes in Young’s double slit experiment?
(a) circular (b) elliptical (c) parabolic (d) hyperbolic
Q.15 The phenomenon which cannot be explained by wave theory of light is
(a) Interference (b) Photoelectric effect
(c) Polarisation (d) Diffraction
Q.16 In which of the following the interference is produced by the division of amplitude?
(a) Lloyd’s mirror (b) Newton’s rings
(c) Young’s double slit experiment (d) Fresnel’s biprism
Q.17 When monochromatic light is replaced by white light in Fresnel’s biprism arrangement, the central
fringe is
(a) dark (b) coloured (c) white (d) none of these
Q.18 In Newton’s rings arrangement, bright and dark rings are obtained using sodium yellow light. What
happens if the top surface of the glass plate on which the lens is kept is highly silvered?
(a) fringes disappear (b) fringe width remains unchanged
(c) fringe width decreases (d) none of these
s HoRt -A nsWeR Q Uestions

Q.1 What is a wavefront?
Q.2 What do you understand by phase difference and path difference?
Q.3 What do you understand by coherent sources?
Q.4 Discuss why two independent sources of light can never be coherent?
Q.5 How are two coherent sources obtained in practice?
Q.6 What are the conditions for interference of light?
Q.7 What are the conditions for maxima and minima in an interference pattern?
Q.8 Define fringe width.
Q.9 Explain fringe width obtained in Newton’s rings experiment.
Q.10 Distinguish between division of wavefront and division of amplitude?
Q.11 What is Fresnel’s biprism?
Q.12 Explain the formation of coherent sources by the use of Fresnel’s biprism.
Q.13 Explain the effect of placing a very thin film in the path of one of the interfering beams.
Q.14 Explain the formation of colours when the white light is incident on a transparent thin film.
Q.15 Why a thick film shows no colours in reflected white light?
Q.16 Explain the necessity of extended source in interference with division of amplitude.
Q.17 What are Newton’s rings?
Q.18 Explain why Newton’s rings are circular.
Q.19 Discuss Michelson’s interferometer.
Q.20 Give the applications of Michelson’s interferometer.
P RACtiCe P RoBleMs
Q.1 Briefly outline the wave theory of light. What is wavefront? How does it propagate?
Q.2 Explain clearly Huygens’ principle for the propagation of light.
Q.3 What are coherent sources? What are the conditions for two sources to be coherent? How are they
realised in practice? Can two independent sources become coherent?
Interference 59
Q.4 Define interference and explain the phenomenon of interference of light.

Q.5 Give the conditions for producing good interference fringes.
Q.6 State and explain in brief the conditions for (a) observance (b) good contrast of fringes and (c) stationary
interference pattern.
Q.7 Discuss briefly the phenomenon of interference with relation to law of conservation of energy.
Q.8 Obtain an expression for fringe width in case of Young’s double slit experiment. Prove that in this case
of interference dark and bright bands are of equal width.
Q.9 A Young’s double slits experiment is carried out with monochromatic light in air. What will be the
change in wavelength and fringe width when the apparatus is immersed in water or the medium is
replaced by an optically denser medium?
Q.10 Discuss how coherent sources are produced with the help of Fresnel’s biprism. Explain with necessary
theory how the wavelength of monochromatic light can be determined using a Fresnel’s biprism?
Q.11 Describe Fresnel’s biprism. Discuss in detail how the wavelength of monochromatic source of light
can be determined with its help?
Q.12 Explain the formation of interference fringes by means of Fresnel’s biprism when a monochromatic
source of light is used, and derive the expression for the fringe width. How will you measure the
wavelength of monochromatic light using biprism method?
Q.13 Describe the geometrical features of Fresnel’s biprism. How can it be used to find the wavelength of
light?
Q.14 Describe the construction, theory and working of Fresnel’s biprism experiment to find the wavelength
of light?
Q.15 Discuss the effect of introducing a thin mica sheet in the path of one of the interfering beams in a
experiment. Deduce an expression for the displacement of the fringes. Describe how this method is
used for finding the thickness of a thin glass plate?
Q.16 Explain the effect of introducing a thin plate of glass in the path of one of the interfering beams in the
experiment. Calculate the displacement of fringes. Show how this method can be used for finding the
refractive index, thickness of the plate and wavelength of light?
Q.17 Using an optical method how would you determine the thickness of a piece of transparent cello tape?
Explain.
Q.18 Discuss the phenomenon of interference in thin films. Obtain the conditions for maxima and minima.
Show that the interference pattern in the reflected and transmitted systems is complimentary. Why an
extremely thin film appears black in reflected light?
Q.19 Why a broad source of light is necessary for observing colours in thin films?
Q.20 Explain analytically the colour of thin films. Why are the colours in reflected and transmitted light
complementary?
Q.21 Describe the interference observed when a thin parallel shaped film is seen by reflected light normally.
Q.22 Explain how interference fringes are formed by a thin wedge shaped film when examined by normally
reflected light. Find the expression for fringe-width. How will you estimate the difference of film
thickness between two points?
Q.23 Explain the formation of Newton’s rings in reflected light?
Q.24 What are Newton’s rings? Explain can these be used to find the wavelength of light?
Q.25 Explain with theory, Newton’s rings method to determine the wavelength of monochromatic light and
discuss which source is preferred, point source or extended source?
Q.26 Discuss the formation of Newton’s rings by reflected light. Describe the experimental arrangement and
give necessary theory. Why are Newton’s rings circular?
Q.27 Explain the phenomenon of interference in thin film and also explain with theory of Newton’s rings
experiment to find the wavelength of monochromatic light.
Q.28 With the help of a neat diagram show an experimental arrangement to produce Newton’s rings by
reflected sodium light. Prove that in reflected light the diameter of the dark rings is proportional to the
square root of the natural number.
Q.29 Describe the principle of construction and working of Michelson’s interferometer.
Q.30 Describe the principle, construction, theory and working of Michelson interferometer to find the
wavelength and the difference in wavelength of a given light.
Q.31 Explain the working of Michelson interferometer. How will you produce circular fringes with it? How
will you measure the difference in wavelength between D lines of sodium light?
Q.32 How do you obtain localised fringes in Michelson’s interferometer? How straight fringes are obtained?
Q.33 How will you find the wavelength of monochromatic light with Michelson’s interferometer?
Q.34 How will you use Michelson’s interferometer to determine the thickness of a thin transparent sheet?
Q.35 Describe the Michelson’s interferometer. How will you use it to calibrate a meter in terms of a standard
wavelength?
U nsolved Q Uestions
Q.1 Calculate the coherence length for a white light when its the wavelength ranges from 4000 Å to 7000
Å. [Ans. 10–15 m]
Q.2 One of the most ideal line of the Krypton has wavelength 6058 Å with line width of 0.0184 Å. Calculate
the coherence length and band width. [Ans. 0.20 m, 1.5 ¥ 109 Hz]
Q.3 The amplitudes of light waves emerging from two slits in Young’s experiment are in the ratio of 1:2.
Find the intensity ratio of the interference patterns. [Ans. 9 : 1]
Q.4 Two coherent sources whose intensity ratio is 81 : 1 produce interference fringes. Deduce the ratio of
maximum to minimum intensity in fringe system. [Ans: 25 : 16]
Q.5 What is the separation between the slits of Young’s double slit experiment that gives second order
maxima at a distance of 5.0 mm from central maxima. The screen is at a distance of 2.0 m from the slit
and the wavelengths of light is 500 nm. [Ans: 0.4 mm].
Q.6 The distance between the slit and biprism and that of between the biprism and screen are each 50 cm.
The obtuse angle of biprism is 179° and its refractive index is 1.5. If the width of the fringes is 0.014
cm, calculate the wavelength of light. [Hint: 2d = 2a (m – 1)] [Ans: 6140Å]
Q.7 Two narrow and parallel slits 0.1 cm apart are illuminated with a monochromatic light of wavelength
5893 Å. The interference pattern is observed at a distance of 25 cm from the slits. Calculate the fringe
width. [Ans: 0:147 nm].
–4
Q.8 On introduction of a thin sheet of mica, having thickness 1.2 ¥ 1 0 cm, in the path of one of the
interfering beams in a biprism experiment, the central fringe is shifted through a distance equal to the
spacing between successive bright fringes. Calculate the refractive index of mica (Given l = 6 ¥ 10–7 m).
[Ans. 1.5]
Interference 61
Q.9 In Fresnel’s biprism experiment the fringes of 0.19 mm width are formed on the screen placed at a
distance of 1.0 m from the slits. A convex lens is placed at a distance of 30 cm from the images of
two coherent sources. The separation between the two images was found to be 0.70 cm. Calculate the
wavelength of light used. [Ans. 5550 Å]
Q.10 The distance between two virtual images of a slit formed by a biprism is 0.3 mm. If fringes of width
0.59 mm are formed on a screen placed at a distance of 30 cm from the slit, calculate the wavelength
of the light used. [Ans. 5900 Å]
Q.11 When a glass piece of thickness 3.6 ¥ 10–4 cm is placed in the path of one of the interfering beams in
a biprism experiment, it is found that the central bright fringe shifts through a distance equal to the
width of four fringes. Calculate the refractive index of the piece of glass. Wavelength of light used is
5.4 ¥ 10–5 cm. (Hint: (m – 1) t = nl) [Ans.: 1.6]
Q.12 In a biprisms experiment fringe width is observed as 0.88 mm. What will it become if the distance
between biprism and slit is reduced to 0.82 times its original distance. (Hint: 2D = l D/b =
2a(m – 1)a) [Ans. 1.07 mm]
Q.13 When a thin soap film of refractive index 1.33 is seen by normally reflected light of sodium of
wavelength 5893 Å, it appears to be black. Find the minimum thickness of the film.
[Ans: 2.215 ¥ 10–7 m]
Q.14 A soap film of refractive index 1.43 is illuminated by a white light incident at an angle of 30°. The
refracted light is examined by a spectroscope in which the dark band corresponding to wavelength
6000 Å is observed. Calculate the thickness of the film. [Ans: 2.23 ¥ 10–7 m]
Q.15 A soap film of refractive index 1.33 is illuminated with light of different wavelengths at an angle of
45°. There is a complete darkness for wavelength 5890 Å. Calculate the thickness of the soap film.
[Ans: 2.614 ¥ 10–7m]
Q.16 Using sodium light (l = 5893 Å) interference fringes are formed by reflection from a thin air wedge.
When viewed perpendicularly, 10 fringes are observed in a distance of 1.0 cm. Calculate the angle of
the wedge. [Ans: 2.95 ¥ 10–4 rad]
Q.17 Light of wavelength 6000 Å falls normally on the wedge-shaped film of refractive index 1.4 forming
fringes that are 2.0 mm apart. Calculate the angle of the wedge. [Ans: 1.07 ¥ 10–4 rad]
Q.18 A wedge-shaped air film having an angle of 40≤ is illuminated by monochromatic light and fringes are
observed vertically through a microscope. The distance between two consecutive bright fringes is 0.12
cm. Calculate the wavelength of light used. [Ans: 4656 Å]
th th
Q.19 In Newton’s rings experiment the diameter of 4 and 25 rings are 0.3 cm and 0.8 cm, respectively.
Find the wavelength of light. Given R = 100 cm. [Ans: 5875 Å]
Q.20 In a Newton’s rings experiment fringes are observed in reflected light of wavelength 5.9 m ¥ 10–7 m
The diameter of 10th dark fringe is 0.5cm. Find the radius of curvature of the lens and thickness of the
air film. [Ans. 1.059 m, 2.95 ¥ 10–6 m]
Q.21 In a Newton’s ring experiment the radius of 10th and 20th rings are 0.2 and 0.3 cm, respectively, and the
focal length of the plano-convex lens is 90 cm. Calculate the wavelength of light used in nanometre.
[Ans. 278 nm]
Q.22 In Newton’s ring arrangement a source is emitting two wavelengths l1 = 6 ¥ 10 and l2 = 5.9 ¥ 10–7
–7
m. It is found that nth dark ring due to one wavelength coincides with (n + 1)th dark ring due to the
other. Find the diameter of the nth dark ring if radius of curvature of the lens is 0.9 m.
[Ans. 1.1289 ¥ 102 m]
Q.23 Newton’s ring are formed in reflected light of wavelength 6000 Å with a liquid between the plane and
curved surfaces. If the diameter of sixth bright ring is 3.1 mm and the radius of curvature of the curved
surface is 100 cm, calculate the refractive index of the liquid. [Ans. 1.37]
–7
Q.24 Newton’s rings are formed with reflected light of wavelength 5.895 ¥ 10 m with a liquid between the
plane and the curved surface. The diameter of the 5th dark is 0.3 cm and the radius of curvature of the
curved surface is 1 m. Calculate the refractive index of the liquid. [Ans. 1.31]
Q.25 In Newton’s rings experiment, the diameter of 7th dark ring is 3.4 mm.
(i) Calculate the diameter of 1 6th dark ring.
(ii) If a small amount of liquid is filled between the lens and glass plate, calculate the radius of 7th and
16th bright rings. (Given: mliq = 1.3) [Ans. (i) 5.25 mm; (ii) 3.0 mm, 4.6 mm]
Q.26 Newton’s rings are observed by reflection between two curved surfaces of radii of curvature 1.0 m and
1.5 m, which are in contact in same plane. Calculate the distance between the 8th and 1 8th dark rings
using monochromatic light of wavelength 5890 Å. [Ans: 1.88 mm]
Q.27 A thin film of a material whose refractive index is 1.45 on being introduced in one of the arms of
Michelson’s interferometer causes a shift of 6 fringes. If the wavelength of the light used is 5890 Å,
calculate the thickness of the film. [Ans: 3.926 m]
Q.28 A sheet of CaF2 of refractive index 1.434 is inserted normally is one arm of a Michelson’s interferometer.
At l = 5896 Å, 40 fringes are observed to be displaced. Calculate the thickness of the sheet.
[Ans: 27 m]
Q.29 A thin plate of refractive index 1.4 is introduced is the path of one of the beams of light in Michelson’s
interferometer and it is found that 50 fringes have crossed the line of observation. The wave length of
light used is 5896 Å. Calculate the thickness of the plate. [Ans: 36.85 m]
Q.30 In Michelson’s interferometer 500 fringes cross the field of view when the movable mirror is displaced
through 0.147 mm. Calculate the wavelength of monochromatic light used. [Ans. 5.88 ¥ 10–7 m]
Q.31 Michelson’s interferometer illuminated by light of wavelength 6438 Å is used to measure the distance
between two points. Calculate this distance if 239 fringes cross the reference mark when the mirror is
moved from one point to the other. [Ans: 7.69 ¥ 10–5 m]
Q.32 In a Michelson’s interferometer 790 fringes cross the field of view when the movable mirror is
displaced through 2.33 ¥ 10–4 m. Calculate the wavelength of monochromatic light used.
[Ans. 5898 Å]
Diffraction 2
Learning Objectives
After reading this chapter you will be able to LO6 Demonstrate Fraunhofer diffraction
by a single-slit/double slit/n slits in
LO1 Differentiate between diffraction and diffraction grating
interference
LO7 Analyse the application of diffraction
LO2 Discuss types of diffraction and the grating in determining wavelength of
differences between them light
LO3 Illustrate Fresnel’s Half-period zones LO8 Explain resolving power of an
LO4 Analyse construction, theory and multi- optical instruments – telescope and
focus behaviour of a zone plate microscope through Rayleigh criterion
LO5 Explain Fresnel’s diffraction by circular LO9 Illustrate resolving and dispersive
aperture power of diffraction grating
Introduction
In the previous chapter, it was discussed that in order to produce an interference pattern, superposition
of at least two beams or waves of light is necessary. For obtaining a sustained interference, these waves
should be coherent and therefore, they were developed from a single source and were separated by the
division of wavefront or amplitude. The same effect is implicated in the diffraction of light. The diffraction
of light is described as the clear bending of waves around small obstacle and the spreading of waves to a
certain extent into the region of geometrical shadow when a beam of light passes through a narrow slit.
We can say that the diffraction is any deviation from geometrical optics resulting from the obstruction of a
wavefront of light. Such effects are observed even if the obstacle is not opaque but causes local variations
in the amplitude or phase of the wavefront of the transmitted light. For example, this effect can be seen
when there is a modification in the properties of the medium through which the wave is traveling, like
variation in the refractive index for light waves. Also, small bubbles or imperfections in a glass lens produce
unwanted diffraction patterns when a monochromatic light is transmitted through it. This phenomenon
can be suitably explained only by assuming the wave nature of light. The effects of diffraction are generally
more prominent for the waves when the size of the diffracting object is of the order of the wavelength
of the wave. The diffraction also has negative implications. For example, the edges of optical images are
seen to be blurred by the diffraction. Therefore, the phenomenon of diffraction leads to a basic limitation
in resolution of the instruments like camera, telescope, microscope, etc.
In addition to electromagnetic waves such as visible light, x-rays and radio waves, the diffraction occurs
with other waves also including sound waves and water waves. Water or ocean waves diffract around
jetties and other obstacles. Sound waves can diffract around objects. This is the reason we can still hear
someone calling us even if we are hiding behind a wall.
2.1 YOung’s DOubLe sLit experiment: DiffractiOn Or LO1

interference?
If you recall Young’s double slit experiment, you will notice that in this experiment the light is blocked
everywhere except at the two opening i.e., the wavefront is obstructed. In order to find the resulting
interference fringe pattern, the two openings were treated as point sources. The same is true if long slits are
used whose widths can be considered as points. However, by taking into account the finite size of the slits,
a more complete analysis can be prepared. Then this problem becomes a problem of diffraction, where the
interference pattern gets modified. In order to explain this, we apply Huygens’ principle according to which
every point of a given wavefront of light acts as a source of secondary spherical wavelets. Here, Fresnel added
another assumption that the actual field at any point away from the wavefront can be found by a superposition
of all these wavelets when we take into account both their amplitudes and phases. This is referred to as
Huygens‑Fresnel principle. In view of this, we consider every point of the wavefront emerging from each
slit as a source of wavelets whose superposition generates the resultant field or diffraction pattern at some
point on a screen. Therefore, we consider a continuous array of sources across both the slits for obtaining
the diffraction pattern and do not take the slits as isolated point sources, as done in the case of interference.
2.2 Difference between DiffractiOn anD interference

In simple words, the diffraction is the bending of light around an obstacle, whereas the interference is the meeting
of two waves. In the phenomenon of diffraction, the interfering beams originate from a continuous distribution
of sources as is clear from Huygens’ principle. However, the interfering beams originate from a discrete number
of sources in the phenomenon of interference. Interference pattern is obtained by the superposition of waves
coming from two different wavefronts originating from the same source. However, the waves emerging from
different parts of the same wavefront superimpose with each other to produce the diffraction pattern. The widths
of the diffraction fringes are not equal, but the widths of the interference fringes may or may not be equal. If you
focus on the points of minimum intensity, you will observe them as perfectly dark in the interference, but these
points in the case of diffraction are not perfectly dark. Moreover, the bright fringes in the interference pattern
are of uniform intensity but these are not of the same intensity in the diffraction pattern.
2.3 tYpes Of DiffractiOn LO2

In order to obtain the diffraction pattern on a screen, we need a source of light, obstacle or aperture and
the observation screen. Here, it is obvious that the distance of the source and the screen from the aperture
will determine the diffraction pattern. Depending upon the distance of the source from the aperture, the
Diffraction 65
wavefronts will reach the aperture either in the spherical form or in the plane form. The same is applicable to
the wavefronts reaching the observation screen after emerging from the aperture. Based on these distances and
hence the shapes of the wavefronts, the diffraction pattern is classified into two classes, namely Fraunhofer
diffraction and Fresnel diffraction.
2.3.1 fraunhofer Diffraction
We need to use plane wavefronts in order to obtain this type of diffraction. This is possible if both the source
of light and the screen are effectively far enough from the aperture so that the wavefronts reaching the aperture
and the observation screen can be considered plane. Then the source and the screen are said to be at infinite
distances from the aperture. This condition can also be attained by using two convex lenses, out of which one
makes the light from the source parallel before it falls on the aperture and the other helps focusing light after
diffraction on the observation screen. This is clear that under the said arrangement the incident wavefront is
plane and the secondary wavelets originating from the unblocked portions of the wavefront are in the same
phase at every point in the plane of the aperture. Here, the diffraction is produced by the interference between
parallel rays that are focused with the help of convex lens. This Fraunhofer type of diffraction is also called
far-field diffraction and it is encountered in the case of a plane transmission grating (discussed later) or
concave reflection grating.
P
O1
a
S O
C
Q
Screen
Figure 2.1
2.3.2 fresnel Diffraction
If the source of light or the observation screen or both of them are at finite distances from the diffracting
aperture, then the wavefronts falling on the aperture or reaching the screen will not be plane. These will
be either spherical or cylindrical depending upon the situation. The diffraction obtained under this type of
arrangement is called Fresnel diffraction for which the curvature of the wavefronts is important. This Fresnel
type of diffraction is also called near-field diffraction. In this arrangement, lenses are not used to make the
rays parallel or convergent. Therefore, the phase of secondary wavelets is not the same at all points in the plane
of the diffracting aperture. Here the resulting field or the diffraction pattern is obtained by the superposition
of these secondary wavelets emanating from different elements of unblocked portions of the wavefront. This
was fantastically explained by Fresnel based on some assumptions. For example, he considered division of
a wavefront into a large number of small area elements or zones called Fresnel zones. Under this situation,
at any point O1 on the screen (Figure 2.1), the resultant field will depend on the combined effect of all the
secondary waves emanating from these zones. The effect of a particular zone at any point will depend on the
distance between the point and the zone. Finally he considered the obliquity (angle a in the figure) of the
point O1 and took the obliquity factor as proportional to (1 + cos a). Therefore, for an elementary wavefront
at C, the resulting effect is maximum at O where a = 0° and cos a = 1. This effect becomes less significant
when we move towards O1, as the angle a gets increased. In the direction tangential to the wavefront (dashed
line in the figure), this effect is one half of that at O, as the angle a = 90°. Based on the obliquity factor,
Fresnel could explain the non‑existence of the wave in the backward direction, where a = 180° and hence the
resulting effect is zero as 1 + cos a = 0
2.3.3 Difference between Fresnel and Fraunhofer Diffractions

Fraunhofer diffraction is encountered in the case of gratings that contain number of slits. Fresnel type of
diffraction is obtained when light suffers diffraction at a straight edge, a thin wire, a narrow slit, a small hole
etc. These effects are the simplest to observe experimentally. For this reason Fresnel diffraction was historically
the first type to be investigated. However, its explanation involves much more difficult mathematical theory
than that required in handling the plane waves of Fraunhofer diffraction. The other difference between the
Fraunhofer and Fresnel diffractions is related to the modification in the pattern when the observation screen
is moved relative to the aperture. In the far‑field approximation (Fraunhofer diffraction), when the screen
is moved, the size of the diffraction pattern scales uniformly and the shape of the pattern does not change.
However, in the near‑field approximation (Fresnel diffraction), both the size and shape of the pattern depend
on the distance between the diffracting aperture and the observation screen.
2.4 fresneL’s HalF-perioD Zones LO3

The concept of half‑period zones was given by Fresnel, based on which he could find the effect produced by
a slightly divergent spherical wave or wavefront at a point ahead of the wave, for example on a screen. Since
every point on the wavefront is treated as the origin of secondary wavelets, the method of half‑period zones
is used to calculate the effect of all these wavelets at a point O (Figure 2. 2). For this we divide the wavefront
into a large number of zones, called Fresnel ’s half‑period zones. For constructing these zones, we consider
Q
Nn v + nl /2
v + 3l /2
P N3 v + 2l /2
v + l /2
N2
N1
C v O
S
Figure 2.2
Diffraction 67
PQRS as a spherical wavefront of a monochromatic light of wavelength l traveling toward the screen. We
draw a perpendicular from the point O to the wavefront. This meets the wavefront at a point C at a distance
of v. We draw a series of circles around C such that their distances from C are CN1, CN2, CN3, ..., CNn and
each circle is a half wavelength farther from O. Therefore, the circles will be at distances v + l/2, v + 2l/2,
v + 3l/2,..., v + nl/2 from O. This way the areas of the zones, i.e., the areas of the rings between successive
circles, are equal (proved later). The area enclosed between CN1, N1N2, N2N3 etc. are called first, second, third
half‑period zones, etc. respectively. Actually the difference of half a period in the vibrations from successive
zones is the origin of the name half‑period zones.
2.4.1 radii and areas of Half-period Zones

In most of the cases the wavelength of light l remains much smaller compared with the distance v. Under
this condition i.e., when l << v, we can calculate the radii and the areas of the zones. The radius of the first
half‑period zone
CN1 ª ON12 - OC 2 = (v + l /2) 2 - v 2 ª vl
Here we have neglected a term l2/4 in view of the condition l << v. Similarly the radius of the second half‑
period zone
CN 2 ª ON 22 - OC 2 = (v + 2l /2) 2 - v 2 ª 2vl
Thus the radius of nth half‑period zone would be

CN n ª nvl
Based on the relations for the radii of the half‑period zones, we can find the area of first, second, third half‑
period zone, etc. Since the radius of first half‑period zone is CN1, its area would be
a1 = p CN12 = p vl
Similarly the area of second half‑period zone is obtained as
a2 = p (CN 22 - CN12 ) = p vl
This is clear from the above that the areas of all zones are the same under the approximation l << v and are
independent of n. However, a more exact evaluation shows that the area gets increased very slowly with n.
2.4.2 resultant amplitude due to Whole Wavefront

As per Huygens’ principle every point of the wavefront sends secondary wavelets in the same phase. However,
these will not reach the point O in the same phase as they travel different distances. Since the average distance
of O from any two consecutive zones differ by l/2 so, the successive zones will produce resultants at O
differing by a phase difference of p. If An be the resultant amplitude of the light due to nth half‑period zone,
then its successive values will have alternate signs because of a phase change of p (or the reverse direction of
amplitude vector) in moving from one zone to the next. It means if the amplitudes of the waves emerging from
the first, second, third half‑period zones, etc. are respectively A1, A2, A3, etc., then the resultant amplitude A
due to the whole wavefront would be
A = A1 – A2 + A3 – ... – An (if n is even)
A = A1 – A2 + A3 – ... + An (if n is odd)
or A = A1 – A2 + A3 – ... + (–1)n – 1 An
Now as the magnitude of successive amplitudes goes on decreasing with the higher order of zones due
to the increased average distance of the zone from O and the larger obliquity, the amplitude A2 is slightly
smaller than A1 but slightly greater than A3. Therefore, to the first approximation, we can assume
A + A3 A + A5
A2 = 1 , A4 = 3 , etc. In view of this, we can expand the above series as
2 2
A Ê A + A2 ˆ Ê A + A5 ˆ An -1
A = 1 + Á - A2 + 1 ˜ + Á - A4 + 3 ˜ +º+ - An (if n is even)
2 Ë 2 ¯ Ë 2 ¯ 2
A Ê A + A3 ˆ Ê A + A5 ˆ A
A = 1 + Á - A2 + 1 ˜ + Á - A4 + 3 ˜ + º + n (if n is odd)
2 Ë 2 ¯ Ë 2 ¯ 2
If the number n is sufficiently large, then the effect due to nth zone would become insignificant and the
A
resultant amplitude due to the whole wavefront can be approximated as A = 1 and hence the intensity
2
A2
I = 1 . Thus the amplitude due to whole wavefront at point O is just half of the one due to the first half‑
4
period and the intensity is equal to one forth of the intensity due to the first half‑period zone a the point O.
2.5 Zone plate LO4

A zone plate is an optical device that is used to verify the correctness of Fresnel’s method of dividing a
wavefront into half‑period zones. It is a special diffracting screen designed to obstruct the light from the
alternate half‑period zones. It is a transparent plate on which a series of concentric circles are drawn with
their radii proportional to the square root of natural numbers. This way the formed alternate annular zones
are blocked. This type of plate behaves like a convex lens and produces an image of a source of light on the
screen placed at a suitable distance.
2.5.1 Construction of Zone plate

On a sheet of white paper we draw concentric circles with their radii proportional to the square root of natural
numbers. The alternate zones are painted black and a significantly reduced photograph of this drawing is
obtained on a glass plate. The negative, when kept in the light path from a distant point source, produces a
large intensity at a point on its axis at a distance determined by the
size of the zone and the wavelength of light used. We can construct
two types of zone plates, namely positive zone plate and negative
zone plate. If the odd zones are transparent and even zones are
opaque on the zone plate, the plate is called positive zone plate. If
the even zones are transparent and odd zones are opaque on the zone
plate, the plate is called negative zone plate as shown in Fig. 2.3(A)
(i) and (ii)). (i) (ii)
Figure 2.3A
2.5.2 theory of Zone plate
Let S be a point source of monochromatic light sending out spherical waves of wavelength l and O is the
position of screen for a bright image (Figure 2.3(B)). Consider an imaginary transparent plate PQ at a right
angle to the line SO passing through the point C. Our aim is to determine the intensity of light at O due to the
wavefront emerging from the source S. In order to find this, we divide the wavefront into half‑period zones
bounded by circles having centres at C and radii equal to CN1, CN2, CN3, ..., CN n as r1, r2 , r3 ,..., rn. These
radii divide the plate PQ into half‑period zones such that from one zone to the next there is an increasing path
difference of l/2. So
Diffraction 69
SN1 + N1O = SC + CO + l/2 (Q l is small)

SN2 + N2O = SC + CO + 2l/2, …,
SNn + NnO = SC + CO + nl/2.
P
Nn
N3
N2
N1
S O
C
k u v
Figure 2.3B
To find the radius rn of the nth circle or the nth zone, we have
SNn + NnO = SC + CO + nl/2 (i)
If we take SC = u and CO= v, then
SNn = ÷(SC2 + CN2) = ÷(u2 + rn2)
1/2
Ê r2 ˆ rn2
= u Á1 + n ˜ ªu+ (Q rn << u )
Ë u2 ¯ 2u
Similarly,
rn2
N nO ª v +
2v
Substituting these values in Eq. (i), we get
rn2 r2 nl
u+ + v + n = u+ v +
2u 2v 2
1 1 nl
or + =
u v rn2
After applying the sign convention, the above formula takes the form
1 1 nl
+ = (ii)
v u rn
nluv
or rn2 = (iii)
c+v
Since u, v and l are constants for a given light, object and image
rn µ n (iv)
It is clear from the above relation that the radii of zones are proportional to the square roots of natural
numbers.
Now the area of the nth zone can be calculated as follows
È nluv (n - 1)l uv ˘ pl uv
an = p (rn2 - rn2-1 ) = p Í - = (v)
Îu + v u + v ˙˚ u + v
This relation shows that the area of the nth zone is independent of n. It means for a given object and image
(for u and v to be constants) the area of all the zones remains the same. However, the average distance of the
zone from O and the obliquity increase with the increase in the order of the zone. Therefore, as the order of a
zone increases the amplitude at O due to the zone gets decreased. Recalling the resultant amplitude A due to
the whole wavefront at point O, we get
A = A1 – A2 + A3 – … + (–1)n – 1 An (vi)
Now we focus on the contribution of zone plate where alternate zones, say even zones in case of positive zone
plate, are blocked. Then the resultant amplitude at O would be
A = A1 + A3 + A5 + … (vii)
Here it is clear that the resultant amplitude A is positive. However, A will be negative if odd zones are blocked
(negative zone plate). Based on the sign of the resultant amplitude A, the zone plates were named as positive
and negative zone plates. A comparison of Eq. (vii) with equation (vi) shows that the resultant amplitude
produced by a zone plate (where the light is blocked by alternate half‑period zones) is greater than that
due to wholly unobstructed wavefront. Hence, the intensity at O is very much enhanced, i.e., the point O is
extremely bright and can be said to be the image of S. This concentration of light at an axial point shows that
the zone plate operates as a lens with O as a focal point. This explains the focusing action of the zone plate.
In order to find the focal length of the zone plate, we concentrate on Eq. (ii) and observe that it is similar to
the lens formula
1 1 1
+ = (viii)
v u f
So a comparison of Eq. (viii) with Eq. (ii) gives
rn2
fn = (ix)
nl
This expression determines the focal length of a zone plate. Since the wavelength l appears in the above
expression, a zone plate will have severe chromatic aberrations. So, it behaves like a convergent lens.
2.5.3 Multi-focus Behaviour of a Zone plate

For an object of infinity, i.e., when u = •, the radius of nth zone can be obtained from Eq. (iii) as
rn = nl v (x)
A comparison of Eq. (x) and Eq. (ix) infers that the distance v is the first focal length f1 (when n = 1), given by
r12
f1 =
l
Diffraction 71
From Eq. (x) it is clear that for fixed rn, the number n gets increased if we reduce the distance v. It means as
the field point O is brought towards the zone plate along the axis, the same zonal area of radius r1 will include
more half‑period zones. If the field point O is brought at a distance of f1/2, n = 2 satisfied the relation (x) for
the same zonal radius r1. Therefore, each of the original zone in this case will now contain two half‑period
zones. For each original zone these two half‑period zones contribute light at the focal point v = f1/2 out of
phase by p with each other. So they cancel and no light is focused by the zone plate at this focal point f1/2.
If we keep on moving the field point O towards the zone plate, we will find n = 3 when v = f1/3 for the same
zonal radius r1. In this case now three half‑period zones are contained in each of the original zones. Out of
these three zones, the effect of two will be canceled due to a phase difference of p between them and the light
will be focused at point O(v = f1/3) due to only one half‑period zone. For the further movement of point O, we
will find no light at v = f1/4, light at v = f1/5, etc. Therefore, we can conclude that a zone plate has multiple
foci of focal lengths f1, f1/3, f1/5, etc. For v = f1/3 the contribution of each original zone is subdivided into
three half‑period zones at the observation point O. So the resultant amplitude will be
A = (A1 – A2 + A3) + (– A4 + A5 – A6) + (A7 – A8 + A9) – … (xi)
In the above expression the first parenthesis is due to the first zone, second parenthesis is due to the second
zone, third parenthesis is due to the third zone, etc. For the zones that are reproduced on a smaller scale,
the obliquity factor is not very important and we may estimate Aj = A1 where j = 2, 3, 4, … Therefore, the
resultant amplitude at v = f1/3 would be simply equal to A1. However, at v = f1 the amplitude will be three
times of this amplitude. Thus, the amplitude at v = f1/3, zone by zone, is reduced by a factor of 1/3 and hence,
the intensity at this point is 1/9 that at v = f1. This can be extended to point at f1/5 also, where the original zone
of radius r1 will include five half‑period zones. Thus the maximum intensity points along the axis and hence
r2
the foci of the zone plate can be found at f n = 1 with n as odd number.
nl
2.5.4 Comparison between Zone plate and Convex lens
Based on the above discussion an understanding is developed that a zone plate operates as a convex lens.
A zone plate has some similarities as well as some differences with a convex lens. For example, both show
chromatic aberration as their focal lengths depend upon the wavelength l. Also, the relation between the
conjugate distances is similar for both of them. However, the differences between these two are listed below.
(i) In case of a zone plate the image is formed by diffraction whereas the rays in case of a lens are
brought to focus by refraction.
(ii) The image due to a convex lens is more intense than due to a zone plate.
(iii) A convex lens has only one focus, whereas a zone plate has n number of foci of reduced intensity
between the points O and C.
1 Ê 1 1ˆ
(iv) The focal length of a lens is given by the relation = ( m - 1) Á - ˜ , where m is the refractive
f Ë R1 R2 ¯
index of the material of lens and R1 and R2 are the radius of curvatures. However, the focal length of
1 nl
the zone plate is given by = 2
f rn
1 v
(v) The focal length of a lens is directly proportional to the wavelength l as f µ fi f µ fi f µ l.
m c
However, the focal length of a zone plate is inversely proportional to l.
(vi) Light takes the same time to go from S to O when passed through any part of the lens. However, in
a zone plate, light from any transparent zone reaches the point O one period later than the light from
the next inner zone.
2.5.5 phase reversal Zone plate

If instead of blocking the alternate zones of the zone plate we introduce an additional path of l/2 between the
waves emanating from successive zones, the amplitudes from these successive zones will be in phase with
each other. Thus we will get four times more intense image than an ordinary zone plate. Such type of zone
plate is called phase reversal zone plate and was investigated by Wood in 1898.
2.6 fresneL’s DiffractiOn bY a circuLar aperture LO5

Let us consider Fig. 2.4 where the aperture AB is taken as circular. In order to study the diffraction through this
circular aperture, Fresnel devised a method for calculating the contribution of various parts of the wavefront
(arc PQ). He divided the aperture into zones with circular symmetry about the axis SCO along with C as the
pole. As discussed earlier, the zones are nothing but the circles on the wavefront separated in such a way that
each zone is l/2 farther from the field point O than the previous zone. In the Fig. 2.4 XY is the screen at which
the diffraction pattern is obtained and the shadow is seen above M and below N on the screen.
If we choose the size of the aperture AB such that it allows only the first half‑period zone, then the resultant
amplitude at point O will be A1. However, if the whole wavefront is exposed and there are a large number of
half‑period zones, then the resultant amplitude at O will be
A = A1 – A2 + A3 – A4 + … = A1/2 (i)
In case of even and odd numbers of half‑period zones, however, the resultant amplitude will be
A1 An -1 A A
A= + - An (if n is even), A = 1 + n (if n is odd) (ii)
2 2 2 2
Therefore, we can conclude that if n is large as in the case of large aperture, An approaches to zero for both
either n is odd or even and the resultant amplitude is half of that the first contributing zone.
M
P
A
S O
C
B
Q
N
Figure 2.4
Diffraction 73
C O
S
xa xb
Y
Figure 2.5
In order to find out the intensity at the point O due to the whole wavefront, we take help of Fig. 2.5. Here the
distances SC = xa and CO = xb and the radius of the circular aperture AB is CA = CB = r. We can find the path
difference between the waves reaching at O as follows
d = SA + AO - ( SC + CO)
d = ( xa2 + r 2 ) + ( xb2 + r 2 ) - xa - xb
1/2 1/2
Ê r2 ˆ Ê r2 ˆ
d = xa Á 1 + 2 ˜ + xb Á1 + 2 ˜ - xa - xb
Ë xa ¯ Ë xb ¯
Ê r2 ˆ Ê r2 ˆ
d = xa Á 1 + 2 ˜ + xb Á 1 + 2 ˜ - xa - xb
Ë 2 xa ¯ Ë 2 xb ¯
r2 Ê 1 1ˆ
d= Á + ˜ (iii)
2Ë a x xb¯
If the position of screen is fixed and the size of the aperture is such that it contains n number of half‑period
zones, then the path difference will be
nl (iv)
d=
2
From Eqs. (iii) and (iv), we obtain
nl r 2 Ê 1 1ˆ
= Á + ˜ (v)
2 2 Ë xa xb ¯
r2 Ê 1 1ˆ
n= Á + ˜ (vi)
l Ë xa xb ¯
Now the number n can be calculated with the help of Eq. (vi) and then with the help of it the resultant
amplitude at the point O can be evaluated. Moreover, from Eq. (v) the position of screen where the intensity
would be either maximum or minimum can be obtained as
xa r 2
xb = (vii)
nl xa - r 2
If the number n is odd, then as per characteristic feature of the half‑period zones the corresponding value of
xa will give the position of the screen such that the point O is bright. If n is even, then the corresponding value
xa will give the position of the screen such that the point O is dark.
1 1 nl
From Eq. (v) we obtain + = . Since this resembles the lens formula, we can find the first focal
xa xb r 2
r2
length as f1 = (for n = 1 and when the source S is at infinity, i.e., when xa = •) and the radius r = nl xb
l
. Now we can analyze the situation when we move the screen towards the aperture. For the fixed width of
the aperture (i.e., the value r) the number of half‑period zones within the aperture will alternately be even
and odd as per relation r = nl xb . Therefore, the point O will alternately be dark (for even values of n)
and bright (for odd values of n). However, if for some distance of the screen (value xb) the aperture contains
only a fraction of the first zone, then the light will spread to the geometrical shadow of the aperture. On the
other hand, now we fix up the position of the screen and change the width of the aperture. When a circular
aperture has the diameter same as the first half‑period zone, then the resultant amplitude at the point O will
be equal to A1. If we make the aperture wider such that it has two half‑period zones, there will be almost
zero amplitude at the point O. Now if we remove all the opaque shields so that all zones of an unobstructed
wavefront contribute, the resultant amplitude will become A1/2 and hence the intensity one fourth of that
due to the first zone aperture alone. These are some curious result because they are not so clear in ordinary
experience. Another result of historic interest is achieved when a round obstacle or the disc is substituted such
that it covers only the first half‑period zone. So all the zones except the first will contribute to the resultant
amplitude at O. As now the second zone is the first contributing zone, the light of amplitude A2/2 focuses at
the point O. Therefore, the intensity at the centre of the shadow of the obstacle will be almost the same as
without disc!
2.6.1 intensity at nonaxial point

In order to calculate the intensity at a nonaxial point such that O1 on the screen XY at a distance of y from the
central point O, we consider the circular aperture AB of radius r (Fig. 2.6). The screen is at a distance of xb.
We consider the source at an infinite distance from the aperture so that the incident wavefront is plane. As
discussed earlier, the pole of the wavefront is defined as the point where the perpendicular from the point O
meets the wavefront. So for the observation point O, the pole is C. However, in moving from point O to O1,
it can be noted from the figure that the pole of the wavefront C is shifted to the position C1. In this situation,
the aperture will not be symmetrical with respect to the half‑period zones. This can be viewed as if the whole
system of half‑period zones is lifted up by the distance y as CC1 = OO1 = y. Thus the upper portions of the
wavefront are now cut off by the obstacle between the points A and C1. Obviously the distance AC1 will
determine the number of zones that are cut off and an increase in this distance shall lead more number of half‑
period zones to cut (Fig. 2.7). Under such situation, in addition to some fully exposed lower order half‑period
zones, the lower portion of higher order zones will contribute to the intensity at points O and O1.
Diffraction 75
O2
4
C1 O1 3
2
A y 1
C O
xb
5
6
Y
Figure 2.6 Figure 2.7
Let us find the maximum intensity at point O along the axis at a particular distance of the screen from the
aperture (Fig. 2.6). Under this condition, let the aperture contain odd number of half‑period zones (say five).
As we move away from the axis from O to O1 (nonaxial point), the pole C gets shifted to C1. Suppose now
first four half‑period zones are fully exposed together with about half of the fifth and sixth zones. Then, the
resultant amplitude at point O1 would be
A5 A6 A1 Ê A1 + A3 ˆ Ê A + A5 ˆ A
A = A1 - A2 + A3 - A4 + - = +Á - A2 ˜ + Á 3 - A4 ˜ - 6
2 2 2 Ë 2 ¯ Ë 2 ¯ 2
A1 A6
or A= - (viii)
2 2
Therefore, the point O1 will have minimum intensity. Now we move further to point O2 so that the first three
half‑period zones are fully exposed together with nearly half of the forth, fifth, sixth and seventh zones.
Then the resultant amplitude will be maximum. Therefore, we obtain series of points along OX at which the
intensity is alternately maximum and minimum. The same is true when we move toward Y, i.e., along OY.
Therefore, we finally observe bright and dark rings of unequal widths about the point of observation O. If the
aperture is large, these rings are seen only near the limits of the geometrical shadow and the intensity falls
off rapidly within the shadow. However, in case of a small aperture which contains only a fraction of the first
half‑period zone, the observation point O will be bright and no rings within the geometrical image of the
aperture will be seen.
The positions of the bright and dark rings can be obtained with the help of Fig. 2.6. Here the path difference
between the secondary waves diffracted from A and B and focused on the point O1 is
2
d = BO1 – AO1 = [(y + r)2 + xb2 ]1/2 – [(y –r)2 + xb ]1/2
È ( y + r )2 ˘ È ( y - r )2 ˘
d = xb Í1 + ˙ - xb 1+
Í ˙ (as the distance y << xb)
Î 2 xb2 ˚ Î 2 xb2 ˚
2ry
d= (ix)
xb
If the path difference d is an even multiple of l/2, the intensity at point O1 will be minimum as the even
2ry l
number of half‑period zones result into almost zero intensity. This condition is satisfied when = 2n .
xb 2
However, the condition of maximum intensity will be satisfied if the path difference is an odd multiple of l/2
2ry l
(odd number of half‑period zones), i.e., when = (2n + 1) . In view of this, the radii of the bright and
dark rings are obtained as xb 2
nl xb
yn = (for dark ring) (x)
2r
(2n + 1)l xb
yn = (for bright ring) (xi)
4r
In terms of the diameter D of the aperture the radius of the first dark ring surrounding the central bright ring
l xb
can be obtained as y1 = .
D
This phenomenon of diffraction at a circular aperture has significance in the formation of images by telescopes
and microscopes. For example, in a telescope the image of a star is seen to consist of bright central disc which
is surrounded by dark and bright rings of gradually diminishing intensities.
2.7 fraunhOfer DiffractiOn bY a singLe sLit LO6

Since in Fraunhofer diffraction the source is effectively at infinite distance, a collimated parallel beam of
monochromatic light of wavelength l can be taken as incident normally on a narrow slit AB of width b (Fig.
2.8). We divide this wavefront into a large number of points in each sending waves of equal amplitude a
according to Huygens’ principle. These waves get diffracted and then interfere to produce diffraction pattern
on the screen. The secondary waves that travel along the direction of incident beam are focused at point C
while those inclined at an angle q with the direction of incident beam (due to diffraction) are focused at
another point P (Fig. 2.8). To find out the resultant intensity at P, we draw a perpendicular AK on BK. It is
clear from the figure that the optical paths of the waves traveled after the plane AK to the point P are equal.
However, the optical paths of the waves originating from points on AB (from A toward B) and reaching the
point P gradually increase. Hence, the phase difference between them gets larger. This is shown in Fig.
2.9(a) for n number of waves each of amplitude a. In this figure a phase difference of f is taken between two
successive waves. The resultant amplitude of these waves at point P is shown as R.
Slit Lens Screen
P
A
x
b C
K
Figure 2.8
Diffraction 77
The total path difference between the waves originating from extreme points A and B is BK = AB sin q = b
sin q. Therefore, the path difference between different waves originating from all the points of the slit AB vary
between zero and b sin q. The phase difference corresponding to path difference b sin q will be (2p/l) b sin q
Since the aperture is divided into n equal parts, the phase difference between any two consecutive parts will
1 2p
be b sin q (= q, say).
n l
P C
a (n 1) f
N a
a
r
f
D 4
R P
a
3f r
C
N
a
2f
B
a a N 2a
a f
O a A O Q
(a) (b)
Figure 2.9
The resultant amplitude and intensity at point P due to all these secondary waves can be obtained by vector
polygon method. Let a be the phase difference between the waves from the initial direction to the resultant
(Fig. 2.9a), then 2a will be the total phase difference between the secondary waves originating from extreme
points of the slit AB (Fig. 2.9b). Here, it is taken that all the amplitudes constitute an arc due to their large
number and small phase difference between them. Because of the symmetry, we have –O = a and –Q = 2a.
The chord OP gives the resultant amplitude due to all the secondary waves at point P.
Then in the DOCN,
ON ON
sin a = =
OC r
or ON = r sin a (i)
where, r is the radius of the circular arc.
\ Chord OP = 2ON = 2r sin a
\ Chord OP = resultant amplitude
\ R = 2r sin a (ii)
The length of the arc ON¢P = na, where n is an integer number and a is the amplitude of each vibration (Fig. 2.9b)
We know that,
Arc ON ¢P na
–PCO = 2a = =
Radius r
na
or 2r = (iii)
a
Substituting the value of 2r in Eq. (ii), we get

sin a
R = na
a
sin a (iv)
R = A0
a
Here we have taken na = A0. Thus, resultant intensity at point P on the screen is given by a measure of
2
sin 2 a È sin a ˘
I = R 2 = A02 = I0 Í
a 2
Î a ˙˚
2
È sin a ˘
I = I0 Í (v)
Î a ˙˚
Thus, the magnitude of the resultant intensity at any point on the screen is a function of a and the slit width
b. Since the phase difference of 2a is introduced due to the path difference of b sin q
2p
\ 2a = ¥ b sin q
l
p
or a= ¥ b sin q
l
2.7.1 Conditions of Maxima and Minima
It is clear from Eq. (iv) that the resultant amplitude R will be a maximum when
sin a
= 1, which s true when a Æ 0, i.e.,
a
p
b sin q Æ 0 or sin q Æ 0
l
fi q = 0°
From Eq. (iv),
A A È a3 a5 a7 ˘
R = 0 sin a = 0 Ía - + - + º˙
a a Î 3 5 5 ˚
È a2 a4 a6 ˘
= A0 Í1 - + + + º˙
Î 3 5 7 ˚
R = A0 for a = 0
Hence the intensity corresponding to a = 0 is I = R2 = A02 = I0. This is called as the intensity of central
(principal) maximum.
In order to determine the position of maximum intensity, let us differentiate Eq (v) w.r.t. to a and equate it
to zero, i.e.,
dI sin a È a cos a - sin a ˘
= 2I0 ˙˚ = 0
da a ÍÎ a2
Since I0 cannot be zero, either sin a = 0 or a cos a – sin a = 0. The equation sin a = 0 determines the positions
of minima [eq. (iv)] except when a = 0 because it corresponds to the position of maximum. Therefore, the
following condition should be satisfied for obtaining the minima
sin a = 0 or = a = ±mp
Diffraction 79
p
or b sin q = ±mp
l
or b sin q = ± ml, m = 1, 2, 3 …
The position of maxima are given by
a cos a – sin a = 0
or a = tan a
This equation can be solved graphically by plotting the curves
y=a
and y = tan a
The first relation y = a represents the equation of straight line passing through the origin making an angle 45°
with the axis and the equation y = tan a represents a discontinuous curve having a number of branches with
asymptotes at the intervals of p (Fig. 2.10). The points of intersections of the these curves will give the values
of a that will satisfy the relation a = tan a.
Therefore, the maxima occur when
3p 5p 7p p
a= , , º or a = (2n + 1) , n = 1, 2, 3, …
2 2 2 2
These are called points of secondary maxima. A measure of intensity of first secondary maxima is obtained
3p
from Eq. (v) with a = , as
2
2
È Ê 3p ˆ ˘
Í sin ÁË 2 ˜¯ ˙ 4
I1 = A02 Í ˙ = 2 I0
Í 3 p ˙ 9p
ÍÎ 2 ˙˚
Similarly, the intensity of second secondary maxima is
2
È Ê 5p ˆ ˘
2
Í sin ÁË 2 ˜¯ ˙ 4
I2 = A Í ˙ = I
2 0
Í 5 p ˙ 25p
ÍÎ 2 ˙˚
2
È Ê 7p ˆ ˘
2
Í sin ÁË 2 ˜¯ ˙ 4
Similarly, I3 = A Í ˙ = I and so on.
2 0
Í 7 p ˙ 49p
ÍÎ 2 ˙˚
Thus, the ratio of relative intensities of successive maxima are
4 4 4
1: : : :º
9p 2 25p 2 49p 2
4
The intensity of the first secondary maxima , i.e., 4.5% that of principal maximum as shown in Fig. 2.11
9p 2
a
y=
Y
an a
y=t
45 p 2p 3p
a
O p 3p 5p 7p
2 2 2 2
Figure 2.10
I
y = an a
y=t
a
a
–3p –2p –p 0 p 2p 3p
Figure 2.11
2.8 fraunhOfer DiffractiOn bY DOubLe sLits

Let a parallel, collimated beam of monochromatic light of wavelength l be incident normally on the two
parallel slits AB and ED (Fig. 2.12) each having a width b and separated by an opaque distance d. The
distance between the corresponding points of the two slits is (b + d). Let the diffracted light be focused by a
convex lens L on the screen XY placed in the focal plane of the lens. As discussed earlier, all the secondary
waves traveling in the direction parallel to MC get focused at point C. Therefore, the point C corresponds to
the position of the central bright maximum.
Diffraction 81
Slits Screen
L X
A P
b/2 q
S1
B
d M C
Eq K
b/2 S q
2
D
Y
Figure 2.12
Further consider that the two slits are equivalent to two coherent sources placed at the middle points S1 and
S2 of the slits AB and ED. Since the resultant amplitude due to a single slit is A sin a/a at any point P making
an angle q with MC, we may consider that each slit is sending a wave of amplitude (A sin a/a). the resultant
amplitude due to interference of these two waves having a phase difference of f¢ at point P can be calculated
as follows. Take S1K as the perpendicular drawn from S1 on S2K. Hence, the path difference between the rays
at point P will be
S2K = (b + d) sin q (i)
The phase difference between them will be
2p 2p
f¢ = ¥ S2 K = (b + d ) sin q (ii)
l l
The resultant amplitude R¢ at point P can be determined by using vector addition method of amplitudes, as
shown in Fig. 2.13.
) B1
(R
de
plitu b
Am Amplitude due
a nt
s ult to slit ED
Re
b f = 2b
C1 A1 E
Amplitude due
to slit AB
Figure 2.13
Since both the slits being of same size send the light of same amplitudes, we may write
sin a
C1A1 = A1B1 = A = R¢
a
Now –B1A1E = f¢ (= 2b, say)
\ Resultant amplitude = C1B1 = R¢ can be obtained as

C1B12 = C1 A12 + A1B12 + 2C1 A1 ◊ A1B1 cos f ¢
2 2
È A sin a ˘ È A sin a ˘ È A sin a ˘ È A sin a ˘
( R ¢)2 = Í ˙ +Í ˙ + 2Í cos f ¢
Î a ˚ Î a ˚ Î a ˙˚ ÍÎ a ˙˚
2
È A sin a ˘
or ( R ¢)2 = 2 Í [1 + cos f ¢ ]
Î a ˙˚
2 A2 sin 2 a
= 2 cos 2 b
a2
4 A2 sin 2 a
= 2
cos 2 b
a (iii)
A measure of intensity can be obtained from the Eq. (iii) as
4 A2 sin 2 a
I = ( R ¢)2 = cos 2 b
a2
sin 2 a
or I = 4I0 cos 2 b (iv)
a2
Based on Eq. (iv), we conclude that the resultant intensity in the pattern depends upon two factors.
sin 2 a
(i) The factor, I 0 , which gives the diffraction pattern due to a single slit.
a2
(ii) The factor cos2 b, which gives the interference pattern in the waves diffracted from the two slits.
From Eq. (iv) it is clear that the maximum intensity I = 4I0, i.e., double slit provides the intensity four times of
the obtained by single slit. Further, it is noted that the intensity I is a product of intensities obtained for double
slit interference and single slit diffractions. Moreover, the expressions of a and b show that the factor cos2a
varies more rapidly than the factor sin2a/a2, as d > b. The product of sine and cosine factors proves that the
double slit diffraction pattern is a modulation of the interference fringe pattern by a single slit diffraction
envelope.
When we analyse diffraction factor sin2a/a2, we find that this gives the principal maximum in the direction
q = 0° on the screen at the point C. This central maximum on its both sides has either alternate minima or
subsidiary maxima of decreasing intensity. The positions of minima are obtained in the direction sin a = 0,
when a π 0. So
a = ±mp
p
or b sin q = ± mp
l
b sin q = ± ml (v)
where m = 1, 2, 3, …, etc. As mentioned m = 0 will give the position of maximum.
As discussed in the case of single slit diffraction the factor sin2 a/a2 gives secondary maxima at the points
3p 5p 7p
a= , , , …. Therefore, the positions of the secondary maxima are obtained in the direction
2 2 2
p
a = (2n + 1)
2
where n = 1, 2, 3, …,
We can analyse the variation of intensity observed by the second factor, cos2 b, as follows.
Diffraction 83
For minimum intensity,

cos2 b = 0 or cos b = 0, i.e.,
p
b = (2m + 1) , m = 0, 1, 2 …
2
Using Eq. (ii), we get
p p
b= (b + d) sin q = (2m + 1)
l 2
p
(b + d) sin q = (2m + 1) (vi)
2
Similarly, bright fringes or maxima are obtained for the values of q for which
cos2 b = 1 or b = ±np
p
or (b + d) sin q = ±np
l
or (b + d) sin q = ±nl (vii)
where n = 0, 1, 2, 3, …
sin2a cos2b
4A2 a2
(c)
a
–3p –2p –p 0 p 2p 3p
(b) cos2b
b
–3p –2p –p 0 p 2p 3p
sin2a
(a) a2
a
–3p –2p –p 0 p 2p 3p
Figure 2.14
sin 2 a
The variations of with a and cos2 b is shown in fig. 2.14. The combined effect of these two, i.e., the
a2
resultant intensity in double slit Fraunhofer diffraction is shown in Fig. 2.14c.
2.8.1 Missing orders in Diffraction pattern

Here the directions of diffraction minima and interference maxima are respectively given by
b sin q = ± ml, m = 1, 2, 3, … (viii)
and (b + d) sin q = ±nl, n = 0, 1, 2, … (ix)
Therefore, the number of interference maxima depends upon the relative values of width of the slits b and
spacing between the two consecutive slits d. If ‘b’ is kept constant and ‘d’ is varied then certain number
of interference maxima are found to be missing from the diffraction pattern. These are known as missing
orders in double slit Fraunhofer diffraction pattern. This happens for a particular set of values of b and d,
which simultaneously satisfy the conditions of diffractions minima [Eq. viii) and interference maxima Eq.
(ix). Since diffraction minima dominate interference maxima the resultant intensity is obtained as zero. The
condition for missing orders is obtained by dividing equation (ix) by equation (viii), i.e.,
b+d n
=
b m (x)
Case-I: d = b, From Eq. (x)
b+d n n
= or = 2
b m m
or n = 2m
Since, m = 1, 2, 3, …, the above condition reads n = 2, 4, 6, … this shows that even order interference maxima
will be missing in the diffraction pattern.
Case-II: d = 2b. From Eq. (x) it follows that
2b + b n n
= or = 3
b m m
or n = 3m
Given m = 1, 2, 3, …, the above condition reads n = 3, 6, 9, …
Thus for d = 2b, 3rd, 6th, 9th, ... order interference maxima will be absent from the diffraction pattern. From the
above discussion it is clear that when the slit separation is an integral multiple of the slit width, the condition
for missing orders is satisfied exactly.
2.9 fraunhOfer DiffractiOn bY N sLits: DiffractiOn grating

We have seen that the intensity produced by double slits is four times of that of single slit. Therefore, it is
expected that the device having more number of slits will produce large intensity. So now we consider a large
number of parallel slits of equal widths separated by equal opaque spaces. Such a device that makes use of
multiple slit diffractions is called the diffraction grating.
It is constructed by ruling large number of fine, equidistant and parallel lines on a optically plane glass plate
with the help of fine diamond point. The ruled lines are opaque to light while the space between any two lines
is transparent to the light and acts as a slit. There are about 15,000 lines per inch in such a grating.
Let a parallel, collimated beam of monochromatic light of wavelength l be incident normally on N‑parallel
slits (grating) each of width b and separated by a opaque distance d. The sum of b and d is known as
Diffraction 85
grating element. The middle points in two consecutive slits separated by the distance (b + d) are known as
corresponding points. Let the diffracted light be focused by a convex lens L on the screen XY placed in the
focal plane of the lens. All the secondary waves traveling in the direction parallel to the direction of incidence
are brought to focus at a point C (Fig. 2.15a). The point C corresponds to the position of central bright
maximum. The rays making an angle q with the direction of incidence are focused at a point P (Fig. 2.15b).
X
P
X S1 q
S2 K1
q
S K2
3 q
C K3
q
C
q
Y SN−1 q
L q L
K N −1 Y
SN
(b)
(a)
Figure 2.15
We may consider that the each slit in the grating is equivalent to an individual coherent source which is placed
A sin a
at the middle of each slit and sending a single wave of amplitude at angle q with the direction of
p a
wave propagation. Here a = ¥ b sin q
l
If S1K1 be the perpendicular on S2K1, then the path difference between the waves originating from S1 and S2
is given by
S2K1 = (b + d) sin q (i)
The corresponding phase difference would be

2p
(b + d) sin q = 2b (say) (ii)
l
It is clear from the above equation that the phase difference between two M
successive waves is constant and is equal to 2b (say). The phases increase in
arithmetical progression (Fig. 2.16). Thus, we can construct the polygon of N
amplitudes. The resultant amplitude and intensity at point P due to the waves C
from N slits can be obtained by vector polygon method and is given by
2b
sin N b A0 sin a sin N b 4b
R = R0 = (iii)
sin b a sin b
A measure of the resultant intensity at point P is given by O
Figure 2.16
A02 sin 2 a sin 2 N b
I = R2 = (iv)
a2 sin 2 b
A02 sin 2 a
The factor gives the intensity distribution in diffraction pattern due to a single slit, while the factor
a2
sin 2 Nb
yields the interference pattern due to N‑slits.
sin 2 b
2.9.1 analysis of Factor sin2 Nb/sin2b showing interference principal Maxima

Now we discuss principal maxima, minima and secondary maxima obtained by diffraction grating.
sin Nb 0
When sin b = 0 or b = ±np, where n = 0, 1, 2, … etc. and sin Nb = 0, we get = . It means it
sin b 0
sin N b
is indeterminate. Therefore, in order to evaluate the value of , we differentiate the numerator and
sin b
denominator according to L’Hospital rule. Thus,
d
(sin N b )
sin N b db
lim = lim
b Æ ± np sin b b Æ ± np d
sin b
db
N cos N b
= lim = ±N
b Æ ± np cos b
sin Nb
By substituting this value of in Eq. (iv), we have
sin b
A02 sin 2 a 2 (v)
I= N
a2
A02 sin 2 a 2
That is the resultant intensity of maxima becomes N . Therefore, the resultant intensity of any of
a2
A2 sin 2 a
the principal maxima in the diffraction pattern can be obtained by multiplying N2 to the factor 0 2 .
a
Being proportional to N2, the brightness of the principal maxima increases with the increase of number of
slits. These maxima are obtained in the direction given by
b = ±np
p
or (b + d) sin q = ±np
l
or (b + d) sin q = ± nl (vi)
where n = 0, 1, 2, …
For n = 0, we get q = 0 which gives zero order principal maximum. For the other values of n as 1, 2, 3, …
we obtain first, second, third, … order principal maximum, respectively. The condition for the existence of
a principal maximum is sometimes called the diffraction grating equation. The value of n gives the order the
diffraction
Minima
The intensity expression (iv) shows that it is minimum when sin Nb = 0 but sin b π 0.
Diffraction 87
Therefore, sin Nb = 0 or Nb = ±mp

or N(p/l) (b + d)sin q = ±mp
or N(b + d) sin q = ±ml (vii)
where m can have all integral values except 0, N, 2N, 3N …, nN, because of these values of m, sin b = 0,
which gives different principal maxima.
It is clear from Eq. (vii) that m = 0 gives principal maximum of zero order. m = 1, 2, 3, … (N – 1) give minima
and m = N gives again principal maximum of first order. Thus, there are (N – 1) equispaced minima between
two consecutive principal maxima.
Secondary Maxima
As discussed there are (N – 1) minima between two successive principal maxima. In order to differentiate two
consecutive minima there should be a maximum between them. Therefore, there would be (N – 2) maxima
between (N – 1) minima. These maxima are known as secondary maxima. The positions of secondary maxima
are obtained by differentiating equation (iv) w.r.t. b and putting it to zero.
dI A2 sin 2 a È sin N b ˘ È N cos N b sin b - sin N b cos b ˘
= 0 2 2Í ˙Í ˙=0
db a Î sin b ˚ Î sin 2 b ˚
or N cos Nb sin b – sin Nb cos b = 0, as sin b = 0 gives principal maxima.
or tan Nb = N tan b
From this relation, we can draw the following triangle.
From Fig. 2.17, we have 1+
N
N tan b
2
tan 2
N tan b b
sin N b =
1 + N 2 tan 2 b
Nb
sin 2 N b N 2 tan 2 b /( 1 + N 2 tan 2 b ) 2 1
Hence, =
sin 2 b sin 2 b Figure 2.17
2
N
=
cos b (1 + N 2 tan 2 b )
2
N2
=
cos 2 b + N 2 sin 2 b
N2
=
1 - sin 2 b + N 2 sin 2 b
sin 2 N b N2
or = (ix)
sin 2 b 1 + ( N 2 - 1)sin 2 b
From Eqs. (iv) and (ix), the intensity of secondary maxima is given by
A02 sin 2 a N2
Is = ◊ (x)
a2 1 + ( N 2 - 1)sin 2 b
N2
It is clear from Eq. (x) that the intensity of secondary maxima is proportional to . Since,
the intensity of principal maxima is proportional to N2. 1 + ( N 2 - 1)sin 2 b
Intensity of secondary maxima 1

=
Intensity of principal maxima 1 + ( N - 1)sin 2 b
2
Hence, as N increases the intensity of secondary maxima decreases. In case of diffraction grating N is very
large. Therefore, the secondary maxima are not visible in the spectrum and there is complete darkness
between two successive principal maxima.
2.9.2 Diffraction pattern
sin 2 a
As mentioned earlier is the diffraction factor and sin2 Nb/sin2 b is the interference factor. In Fig. 2.18
a2
we plot these two separately and also a combined effect (product) of them is shown. Thus the intensity
distribution or the diffraction pattern due to N slits or diffraction grating is shown in Fig. 2.18c.
sin2a
a2
(a)
a
sin2Nb
sin2b
(b)
b
2
A 0 sin2 a sin N b
2 2
a sin2b
(c)
Ist Zero 1st

Order Order Order
Figure 2.18
2.9.3 Missing orders in Diffraction pattern

As observed in double slit diffraction, for a particular set of values of slit width and slit separation, certain
number of interference maxima are found to be missing. In the case of diffraction grating also, similar situation
arises for certain values of b and d. It means the condition of missing order is met in the diffraction pattern
of diffractions grating. For this the conditions of interference maxima and diffraction minima should be
simultaneously satisfied. The conditions of interference maxima is given by equation (vi). For the diffraction
minima, it can be noted from diffraction factor sin2 a/a2 that it should be zero. It means sin a = 0 but a π 0.
Êp ˆ
Therefore, sin Á b sin q ˜ = 0. This would be true if
Ël ¯
Diffraction 89
p
b sin q = ±mp
l
or b sin q = ± ml.
Therefore, in order to meet the condition of missing orders following relations should be satisfied.
(b + d) sin q = nl, n = 0, 1, 2, … (Interference maxima) (xi)
b sin q = ml, m = 1, 2, 3, … (Diffraction minima) (xii)
From Eqs. (xi) and (xii), we get
b+d n
=
b m
Êb + dˆ
n=Á
Ë b ˜¯
or m
This is the condition of missing order of interference maxima in diffraction pattern. The absent orders are
given below.
If d = b, then
Êb + dˆ
n=Á m = 2m
Ë b ˜¯
Therefore, for m = 1, 2, 3, … the 2nd, 4th, 6th, … order interference will be absent as n = 2, 4, 6 … If d = 2b,
then, n = 3, 6, 9, … for m = 1, 2, 3, … Therefore the 3rd, 6th, 9th, … order interference will be absent from the
diffraction pattern.
2.9.4 angular Width of principal Maxima

If qn be the direction of nth principal maxima and (qn + dqn) and (qn – dqn) the directions of first outer and
inner sided minima adjacent to the nth maxima, then the value of angular width will be 2dqn (Fig. 2.19).
M
P2
m = nN + 1
m = nN
dqn
P1
dqn m = nN – 1
qn
O P
Figure 2.19
The directions of nth order principal maxima and minima are given as follows
(b + d) sin qn = nl (xiii)
and N(b + d) sin qn = ml (xiv)
For the first order outer and inner sided minima adjacent to the nth maxima, qn should be replaced with
(qn ± dqn) and m = (nN ± 1). Then from Eq. (xiv), we get
N(b + d) sin (qn ± dqn) = (nN ± 1) l (xv)
or N(b + d) (sin qn cos dqn ± cos qn sin dqn) = (nN ± 1)l (xvi)
For small values of dqn, cos dqn = 1 and sin dqn = dqn. With this equation (xvi) becomes
N(b + d) sin qn ± N(b + d) cos qn dqn = nNl ± l (xvii)
or Nnl ± N(b + d) as on dqn = nNd ± l [with the help of Eqn (xiii)]
With the help of Eq. (xiii), the above equation can be written as
N(b + d) cos qn dqn = l
l
or dq n = (xviii)
N (b + d ) cos q n
2l
or 2dq n = (xix)
N (b + d ) cos q n
This is the expression for angular width of the nth order principal maxima, which shows that it depends on
the total number of lines present on the grating and the wavelength of the light used in addition to the grating
element.
2.10 appLicatiOn Of DiffractiOn grating LO7

In the previous chapter we learnt how to determine the wavelength of monochromatic light using mechanism
of interference. We can also find the wavelength of light by using the concept of diffraction with the help of
diffraction grating.
2.10.1 Determination of Wavelengh of light

The theory of diffraction grating says that the principal maxima for a monochromatic light of wavelength l
is obtained as per the expression (b + d) sin q = nl. It is evident from this expression that if we are able to
accurately measure the angle of diffraction q and the grating element (b + d), which is equal to the reciprocal
of the number of lines per cm, we can determine the wavelength of light. A typical arrangement of using
diffraction grating for the measurement of light is depicted in Fig. 2.20.
In place of a monochromatic light, if we use the ordinary light (which has different wavelengths), then the
beam gets dispersed (diffracted) by the grating. Under this situation the spectrum of constituent wavelengths
is obtained, as shown in Fig. 2.21. This is due to the different angle q for different wavelength l for satisfying
the condition of principal maxima.
Thus by knowing the grating element (b + d), angle of diffraction q and the order n of the diffraction, the
wavelength l of a particular colour can be calculated with the help of formula
(b + d ) sin q
l=
n
Diffraction 91
G pe
s co
Grating Tele
S
q
Collimator
Figure 2.20
Grating
R2
IInd order
V2
R1
Ist order
V1
q2
q1
zero order
q1 (central image)
q2
V1
R1 Ist order
V2
R2 IInd order
Figure 2.21
2.11 resOLving pOwer Of an OpticaL instrument LO8

If two images are very close to each other, they may appear as one and it is impossible for the eye to see them
separately. Therefore, some optical instruments are used to resolve such images. Our eye can resolve two
objects only when the angle subtended by them at the eye is greater than one minute (1/60)°. Here we say that
the resolving limit of the normal eye is (1/60)°. The resolving power of an optical instrument is defined as its
ability to just resolve the images of two close point sources or small object. It is the ability of instrument to
measure the angular separation of two images that are close to each other.
2.11.1 rayleigh Criterion for resolution

According to Rayleigh, two close images or two close spectral lines of equal intensities are said to be resolved
by an optical instrument if the position of the central maxima of one spectral line coincides with the first
minima of the other spectral line and vice‑versa.
The intensity distribution curves (diffraction pattern) of spectral lines having wavelengths l and l + dl of
equal intensities are shown in Fig. 2.22.
l l + dl
(a)
Resultant Intensity Curve
l l + dl
dl
(b)
Resultant Intensity Curve
l
l + dl
(c)
Figure 2.22
Figure 2.22a says that when the difference in the angle of diffraction is large, the two spectral lines can be
seen as separate ones and hence these spectral lines will be well resolved. In Fig. 2.22b, the difference in
the angle of diffraction is such that the principal maxima of one just coincides with the first minima of the
other. Here resultant intensity curve shows a dip in the middle of the central maxima of these spectral lines.
According to Rayleigh, these spectral lines can be distinguished from one another and are said to be just
resolved. If the central maxima of two spectral lines corresponding to the wavelengths l and l+ dl are very
close to each other, as shown in Fig. 2.22c, then these two spectral lines overlap and they cannot be seen as
separate ones.
According to Rayleigh criterion, two images or two close spectral lines of equal intensities are said to be just
resolved when the resultant intensity at the dip is (8/p2)of the intensity of either central maxima. This can be
proved as follows:
Diffraction 93
According to the theory of single slit Fraunhofer diffraction,

sin 2 a
I = I0
a2
Since first minima is formed at an angle a = p, the angle at the point of intersection will be p/2. In this case,
intensity of each curve at the dip
sin 2 (p /2) 4
I1 = I 2 = I 0 = 2 I0
(p /2) 2 p
The resultant intensity at the dip

4 4
I = I1 + I 2 = I + 2 I0
2 0
p p
8
= I0
p2
= 0.81I 0
2.12 resOLving pOwer Of a teLescOpe LO8
Telescope is an optical instrument which is used to produce a magnified image of a distant object. The
telescope can also form separate images of two close small objects situated at large distance. In this context
it is important to investigate its resolving power.
Since a telescope consists of a system of lenses, we consider the diameter of the objective lens AB of the
telescope as a. Further, we take two distant object P and Q such that they substend an angle q on the objective
lens of the telescope as shown in Fig. 2.23.
A
q
q Q0
P P0
q 90°
N
Q
q B
Figure 2.23
Now, we consider that a beam of light is incident on the objective lens of the telescope from these two
neighbouring point sources. The image of each point source gives Fraunhofer diffraction pattern. Let P0
and Q0 be the positions of the central maximum of the images of P and Q. According to Rayleigh criterion,
these two images are said to be resolved if the position of the central maximum of diffraction pattern of
one coincides with first minimum of the other and vice‑versa. All the secondary waves traveling in the
direction AP0 and BP0 will meet on P0 and the path difference between them will be zero. Thus, the point
P0 corresponds to the position of the central maximum of the first image. Similarly, the secondary waves
traveling along AQ0 and BQ0 are met on Q0 and the path difference between AQ0 and BQ0 is equal to BN.
BN = AB sin q = a sin q
or BN = aq (for small angle)
If the path BN = aq be equal to l, then the position of the central maximum of Q0 corresponds to the first
minimum for the first image P0 and hence this condition satisfies the Rayleigh criterion of resolution. Thus
aq = l
or q = l/a (i)
The Eq. (i) holds goods for rectangular aperture.
For circular apertures, the modified from the Eq. (i) is written as
1.22l
q= (ii)
a
where, l is the wavelength of the light and q refers to the limit of resolution and its reciprocal gives the
resolving power of the telescope.
1 a
Resolving power = = (iii)
q 1.22l
The above equation says that the resolving power of a telescope would be higher, if the aperture a of the
objective lens is taken larger.
Limit of Resolution of a Telescope
The limit of resolution is defined as an angle subtended at the objective lens of the telescope by two distant
point objects which are just resolved when seen through the telescope. For the smaller values of this subtended
angle on the objective, the resolving power of the telescope is said to be higher.
2.13 resOLving pOwer Of a micrOscOpe LO8

Microscope is used to produce magnified images of small objects. In simple microscope, a virtual image
is formed by a convex lens of short focal length when an object is placed just inside its principal focus.
However, compound microscope consists of two short focus convex lenses, named the objective and the
eyepiece. The resolving power of a microscope is defined as its ability to form distinctly separate images of
two objects lying close together. In the present case, the objects are not situated at large distance rather the
objects are close to the objective lens. Unlike the case of telescope, the objects here subtend a large angle at
the object plane. In the case of microscope, we mainly try to find out the smallest distance between two point
objects which will produce images that are just resolved.
In Fig. 2.24, P and Q are two point objects separated by a distance d from each other and AB is the aperture of
the objective lens of the microscope. P0 and Q0 are the corresponding Fraunhofer diffraction pattern of these
two images. P0 and Q0 are the positions of the central maximum of P and Q that are surrounded by alternate
dark and bright diffraction minima and maxima. According to Rayleigh criterion, these two images are said
to be resolved if the position of the principal maximum of image P0 corresponds to the first minimum of the
image Q0.
Diffraction 95
Q
d
P P0
N
Q0
B
Figure 2.24
The path difference between the extreme rays originating from the point Q and reaching at P0 is given by
D = (QB + BP0) – (QA + AP0)
= QB – QA [Q BP0 = AP0]
= (QM + MB) – (NA – NQ) (Please refer to Fig. 2.25)
= (QM + PB) – (PA – NQ) [For small distance between P and Q]
= QM + NQ [\ PB = PA]
From the triangles PNQ and PMQ in Fig. 2.25
QM = PQ sin q = d sin q
NQ = PQ sin q = d sin q
Therefore, the path difference = QM + NQ = d sin q + d sin q = 2d sin q.
A
N Q M
qq
C
P
Figure 2.25
According to Airy, if the path difference is equal to 1.22l (for circular aperture), then the maxima of image
P0 coincides with the minimum of the image Q0. Therefore, these two images appear just resolved,. Thus,
2d sin q = 1.22l
1.22l
d= (i)
2 sin q
If the space between the object and objective lens is filled with an oil of refractive index m, then
1.22l
d= (ii)
2 m sin q
Here m sin q is known as numerical aperture (NA) of the objective. Hence,

1.22l
d= (iii)
2 NA
This equation gives the limit of resolution when two objects are self‑luminous.
2.14 resOLving pOwer Of a pLane DiffractiOn grating LO9

Consider a parallel beam of light of wavelengths l and l + dl incident normally on the plane transmission
grating having grating element (b + d) and total number of rulings N (Fig. 2.26). Then the resolving power of
the grating is defined as the ratio of wavelength (l) to the difference dl of the wavelength i.e., l/dl.
The separate diffraction pattern for l and l + dl is shown in the Fig. 2.26. According to Rayleigh criterion,
these spectral lines are just resolved as the principal maxima of one line lies just on the first minima of the other.
Now the direction of nth principal maximum for a wavelength l is given as
(b + d) sin q = nl (i)
th
The direction of n principal maximum for a wavelength l + dl is given by
(b + d) sin (q + dq) = n(l+ dl) (ii)
The equation of minima for wavelength l is
N(b + d) sin q = ml (iii)
X M
l + dl
P2
l P1
dq
q
O P
Y N
Figure 2.26
Here m has all the integral values except 0, N, 2N, … nN, because for these values of m the condition for
maxima is satisfied. Thus, first minimum adjacent to nth principal maximum in the direction (q + dq) can
be obtained by substituting the values of m as (nN + 1) in Eq. (iii). Therefore, first minima in the direction
(q + dq) is given by
N(b +d) sin (q + dq) = (nN + 1)l
Diffraction 97
(nN + 1)l
or (b + d) sin(q+ dq) =
N
l
(b + d) sin (q + dq) = nl + (iv)
N
A comparison of Eq. (iv) with Eq. (ii), i.e., the Rayleigh criterion for just resolution, gives
l
n(l + dl) = nl +
N
l
or nl + ndl = nl +
N
l
or ndl =
N
l
or = nN
dl
This is the required expression for the resolving power of the plane diffraction grating. This says that the
number of lines per cm of a grating should be larger in order to increase its resolving power.
2.15 Dispersive pOwer Of a pLane DiffractiOn grating LO9

As we have seen, light of different wavelengths get dispersed/diffracted by the grating at different angles. In
view of this, the angular dispersive power of a diffraction grating is defined as the rate of change of the angle
dq
of diffraction with the wavelength of light. It is denoted by .
dl
For a plane transmission grating, the condition for principal maxima is
th
where (b + d) is the grating element and q is the angle of diffraction for n order principal maxima.
Differentiation of Eq. (i), w.r.t. l gives
dq
(b + d) cos q =n
dl
dq n
\ fi Dispersive power = = (ii)
d l (b + d ) cos q
Here dq is the angular separation between the two lines having difference dl in the wavelengths.
We can make following observations from Eq. (ii).
(i) The dispersive power is directly proportional to n, i.e., to the order of diffraction. Thus, higher is the
order, greater is the dispersive power. Hence the angular separation of two spectral lines is double in
the second order diffraction as compared to that in the first order.
(ii) The dispersive power is inversely proportional to the grating element (b + d). This means that the
dispersive power is directly proportional to number of lines per cm of grating. Therefore, the angular
dispersive power of two given lines is greater with a grating having larger number of lines per cm.
(iii) The dispersive power is inversely proportional to cos q. thus, if the angle of diffraction q = 0°, cos
q = 1 and hence the angular dispersion is minimum. Therefore, if q is small, the value of cos q may
be taken as unity so the influence of cos q may be neglected.
If we neglect the influence of the factor cos q, then dq µ dl for a given order, i.e., the angular dispersion
of two spectral lines in a particular order is directly proportional to the difference in the wavelengths. Such
spectrum is called a normal spectrum.
Linear Dispersive Power
If dx be the linear separation of two spectral lines differing in wavelengths by dl in the focal plane of a lens
of focal length f, then we have
dx = f dq
dx dq
Here, the linear dispersive power is defined as = f
dl dl
fn
=
(b + d ) cos q (iii)
s UMMARY
The main outcomes of this chapter are summarized as follows:
✦ Initially we discussed the phenomenon of diffraction.
✦ We clarified Young’s double slit experiment with reference to both the phenomena of interference and
diffraction.
✦ A clear distinction of the diffraction from the interference was given.
✦ Depending upon the distance of the source from the aperture, the incident wave can be realized either in
the form of spherical wavefront or plane wavefront. The same is applicable to the wavefronts reaching
the screen after emerging from the aperture. Based on these distances and hence the shapes of the
wavefronts, the diffraction was divided into two classes, namely Fraunhofer diffraction and Fresnel
diffraction.
✦ A concept of finding the resultant of a wavefront on the screen was given by Fresnel in terms of half‑
period zones. So Fresnel’s half‑period zones along with their construction were discussed.
✦ The concept of half‑period zone was extended to zone plate, where alternative half‑period zones are
blocked and intense diffraction pattern is obtained.
✦ It was proved that a zone plate acts like a convex lens but it has multiple foci. Similarities and differences
of a zone plate with a convex lens were summarized.
✦ Theory was given for Fresnel’s diffraction by a circular aperture. Modification in the diffraction pattern
at various nonaxial points on the screen was explained. The diffraction pattern was investigated when
we move the screen toward the aperture by keeping the aperture diameter fixed. All these phenomena
were discussed with the help of concept of half‑period zones.
✦ Fraunhofer diffraction by a single slit and double slits was investigated by deriving the conditions
of maxima and minima. It was discussed how certain number of interference maxima are found to
be absent in the double slit experiment when we play with the slit width and slit separation. These
conditions are called as the conditions for missing orders.
✦ Since the intensity in the case of double slit diffraction was found to be four times of the one obtained
by the single slit, the need of an arrangement having a large number of slits was discussed and the
Diffraction 99
plane diffraction grating was introduced. The concepts of principal maxima and secondary maxima
were brought in.
✦ The resolving power of an optical instrument is defined as its ability to just resolve the images of
two close point sources or small objects. These two close images or two close spectral lines of equal
intensities are said to be just resolved if the position of the central maxima of one spectral line coincides
with the first minima of the other spectral line and vice‑versa. This is called Rayleigh criterion for
resolution.
✦ Expressions for resolving powers of a diffraction grating, telescope and microscope were obtained.
Dispersive powers of a plane diffraction grating, both angular dispersive power and linear dispersive
power, were discussed. Angular dispersion power is defined as the rate of change of the angle of
diffraction with the wavelength of light.
s olved e XAMPles
E xamplE 1 A plane wavefront of light (l = 5000 Å) is incident on an opening and is received on a screen at
a distance of 100 cm from the opening. Find the radius of 80th half‑period zone and the area of a half‑period
zone.
Solution Given l = 5.0 ¥ 10–7 m, v = 100 cm = 1.0 m and n = 80.
Radius of nth half‑period zone
rn = nnl = 80 ¥ 1.0 ¥ 5 ¥ 10-7
= 6.32 ¥ 10–3 m
= 0.632 cm.
Area of half‑period zone
= pvl = 3.14 ¥ 100 ¥ 5 ¥ 10–5
= 0.0157 cm2
E xmaplE 2 Find the radius of the first half‑period zone of a zone plate that behaves like a convex lens of
focal length 60 cm. Given l = 6000 Å.
Solution Given f = 0.60 m, l = 6.0 ¥ 10–7 m and n = 1.
rn2 r12
Formula used is fn = or f1 =
nl l
or r12 = 0.6 ¥ 6.0 ¥ 10–7
r1 = 0.6 ¥ 10–3 m
or r1 = 0.6 mm
E xamplE 3 A parallel beam of light of wavelength 5 ¥ 10–7 m falls on a circular aperture and the diffraction
pattern is observed on a screen 0.30 m away. Find the radius of circular opening so that the intensity of light
on the screen is 4 times the intensity in absence of opening.
Solution Given v = 0.30 m and l = 5 ¥ 10–7 m.
In the given case, the radius of the opening
= radius of 1st half‑period zone
= nl
= 0.30 ¥ 5 ¥ 10-7
= 0.387 ¥ 10–3 m.
= 0.0387 mm
E xamplE 4 Keeping the distance of observation point fixed as 0.50 m, calculate the number of half‑period
zones in a circular opening of radius (a) 2.0 mm (b) 20 mm, where the light of wavelength 6000Å is used.
Solution Given v = 0.50 m, r1 = 2.0 mm, r2 = 20 mm and l = 6.0 ¥ 10–7 m.
Let us consider an as the area of circular opening, which contains n half‑period zones.
an = n(pvl) (Q opening contains n half‑period zones each of which has area = pvl)
an = pr2
For r1 = 2.0 mm = 2.0 ¥ 10–3 m,
p r 2 r12 (2.0 ¥ 10-3 ) 2 4.0 ¥ 10-6 40
n= = = = =
pnl nl 0.5 ¥ 6.0 ¥ 10-7 3.0 ¥ 10-7 3
n = 13
Similarly, for r2 = 20 mm = 0.02 m
r22 (0.02)3
n= =
nl 0.5 ¥ 6.0 ¥ 10-7
4.0 ¥ 10-4
= = 1.333 ¥ 103
3.0 ¥ 10-7
n = 1333
E xamplE 5 The diameter of the first ring of a zone plate is 1.0 mm. If the plane waves of wavelength 5000 Å
fall on the plate, find where a screen should be placed so that light is focused at the brightest spot.
Solution Given r1 = ½ mm = 0.5 ¥ 10–3 m, l = 5.0 ¥ 10–7 m and n = 1.
rn2
Formula used is fn =
nl
Q n=1
r 2 (0.5 ¥ 10-3 ) 2 2.5 ¥ 10-7
\ f1 = 1 = = = 0.5 m
l 5.0 ¥ 10-7 5.0 ¥ 10-7
E xamplE 6 Find the radius of the first three transparent zones of a zone plate behaving like a convex lens of
focal length 1.0 m for light of wavelength l = 5893 Å.
Solution Given f = 1.0 m and l = 5.893 ¥ 10–7 m.
Considering positive zone plate in which odd number of zone are transparent so that
r= fnl
st
Then, for 1 zone, n = 1
\ r1 = 1.0 ¥ 1 ¥ 5.893 ¥ 10-7
= 7.676 ¥ 10–4 m
Diffraction 101
Similarly r3= 1.0 ¥ 3 ¥ 5.893 ¥ 10-7

= 1.329 ¥ 10–3 m.
and r5 = fnl
= 1.0 ¥ 5 ¥ 5.893 ¥ 10-7
= 17.165 ¥ 10–4 m
= 1.716 ¥ 10–3 m.
E xamplE 7 What is the radius of the tenth zone plate of focal length 20 cm for light of wavelength 5000 Å?
Solution Given f = 0.20 m, f = 5.0 ¥ 10–7 m and n = 10.
rn2
Formula used is fn =
nl
or r102 = fnl = 0.2 ¥ 10 ¥ 5.0 ¥ 10–7 m = 10–6
or r10 = 10–3 m
r10 = 1.0 mm.
E xamplE 8 The image of a point source of light (l = 5890 Å) at a distance 1.0 meter from the zone plate is
observed at 2.0 meter on the other side. Calculate
(a) the focal length of the zone plate (b) power of a zero plate
(c) diameter of 1st zone.
Solution Given l = 5.89 ¥ 10–7 m, u = 1.0 m and v = 2.0 m.
(a) Focal length
1 1 1 1 1 3
= + = + =
f u v 1 2 2
2
f = = 0.67 m.
3
(b) Power of zone plate
1
P= [because zone plate acts as a convex lens]
f
1 3
P= = = 1.5 D
2/3 2
(c) Let r1 be the radius of the first zone. Then,
r12 = fnl
2
r12 = ¥ 1 ¥ 5.893 ¥ 10-7 m 2
3
r12 = 3.93 ¥ 107 m 2
r1 = 6.267 ¥ 10-4 m
Hence diameter of 1st zone

= 2r1 = 1.253 ¥ 10–3 m = 1.253 mm
E xamplE 9 If the focal length of the zone plate is 1.0 m for light of wavelength 6.0 ¥ 10–7 m, what will be
the focal length for the wavelength 5.0 ¥ 10–7 m.
Solution Given l1 = 6.0 ¥ 10–7 m and l2 = 5.0 ¥ 10–7 m.
rn2
Formula used is fn=
nl
Thus, focal length f1 and f2 are
rn2 r2
f1 = and f 2 = n
l1 2l 2
or f 2 l1 l1
= or f 2 = f1
f1 l2 l2
6.0 ¥ 10-7
f 2 = 1.0 ¥
5.0 ¥ 10-7
f = 1.2 m
E xamplE 10 An object is placed at 20 cm from a zone plate and the brightest image is situated at 20 cm
from the zone plate with light of wavelength l = 4000 Å. Calculate the number of Fresnel’s zone in a radius
of 1.0 cm of that plate.
Solution Given u = 0.20 m, v = 0.20 m, l = 4.0 ¥ 10–7 m and r = 0.01 m.
rn2 1 1 1
Formula used is fn = and = +
nl f u v
1 1 1 2 0.2
= + = or f = m = 0.1 m
f 0.2 0.2 0.2 2
or f = 10 cm
rn2
So, the number of Fresnel’s zone (n) =
fnl
(0.01) 2
=
(0.1) ¥ 4 ¥ 10-7
1.0 ¥ 10-4 1.0
= = ¥ 104 = 0.25 ¥ 104
4 ¥ 10-8 4
\ n = 2500
E xamplE 11 The diameter of the central zone of a zone plate is 2.3 mm. If a point source of light (l = 5893 Å)
is placed at a distance of 6.0 m from it, calculate the position of the first image.
Solution Given diameter (d) = 2.3 mm = 2.3 ¥ 10–7 m, r = d/2 = 1.15 ¥ 10–3 m, l = 5.893 ¥ 10–7 m and n = 1.
rn2 r 2 (1.15 ¥ 10-3 ) 2
Formula used is fn = = =
nl l 5.893 ¥ 10-7
= 2.244 m
\ f = 2.2 m
1 1 1 1 1
= + = +
f u v 6 v
Diffraction 103
1 1 1 6 - 2.2 3.8
= - = =
v 2.2 6 2.2 ¥ 6 13.2
13.2
v= = 3.47
3.8
v = 3.47 m
Hence, first image is formed at a distance of 3.47 m.
E xamplE 12 A zone plate is made by arranging the radii of the circles which define the zones such that they
are the same as the radii of Newton’s rings formed between a plane surface and the surface having radius of
curvature 200 cm. Find the principal focal length of the zone plane.
Solution Given R = 2.0 m.
r12
Formula used is f1 = (Q n = 1) (i)
l
By Newton’s ring formula
rn = nl R
For n = 1
r1 = l R (ii)
By using Eqs. (i) and (ii), we have

lR
f1 = = R = 2.0 m
l
f1 = 2.0 m
E xamplE 13 A single slit of width 1 mm is illuminated by light of wavelength 589 nm. Find the angular
spread of the central maxima of diffraction pattern observed.
Solution Given l = 5.89 ¥ 10–7 m and slit‑width (b) 1.0 ¥ 10–3 m.
Formula used is b sin q = ml
For first minima, m=1
\ sin q = l/b
È 5.89 ¥ 10-7 ˘
q = sin–1 [l/b] = sin–1 Í -3 ˙
Î 1.0 ¥ 10 ˚
u = 0.03374°
The angular spread of central maximum is 2q
\ 2q = 2 ¥ 0.03374
2q = 0.0675°
Therefore, angular spread of central maximum is 0.0675°.
E xamplE 14 In Fraunhofer type diffraction at narrow slit of width 0.2 mm, a screen is placed 1.2 m away
from the slit. In the fringe pattern, first minimum lie at 3.7 mm on either side of the central maximum. Find
out the wavelength of light.
Solution Given b = 2 ¥ 10–4 m, D = 1.2 m and x = 3.7 ¥ 10–3 m.
A q
x
q
q
Figure 2.27
From the given figure, if q is very small, then

x
sin q = (i)
D
condition of minima b sin q = ml
For m = 1, then sin q = l/b (ii)
By using Eqs. (i) and (ii), we get
l x xb 3.7 ¥ 10-3 ¥ 2 ¥ 10-4
= or l = =
b D D 1.2
l = 6167 Å
E xamplE 15 A light of wavelength 550 nm falls normally on a slit of width 2.2 mm. Determine the angular
position of second and third minima.
Solution Given b = 2.2 ¥ 10–6 m and l = 5.5 ¥ 10–7 m.
The formula used is b sin q = ml
For angular position of second minima
2l 2 ¥ 5.5 ¥ 10-7
sin q 2 = = = 0.5
b 2.2 ¥ 10-6
or q2 = sin–1 (0.5) = 30°
or q2= 30°
Similarly, for angular position of third minima
3l 3 ¥ 5.5 ¥ 10-7
sin q 3 = = = 0.75
b 2.2 ¥ 10-6
u3 = 48.59°
E xamplE 16 In Fraunhofer diffraction at a slit of width 1.2 ¥ 10–6 m, find the half‑angular width of the
central bright maximum if the slit is illuminated by light of wavelength 5890 Å.
Diffraction 105
For first minima m = 1

5.89 ¥ 10-7
sin q = l / b = = 0.491
1.2 ¥ 10-6
q = sin–1 (0.491)
or q = 29.41°
E xamplE 17 Parallel beam of light (5000 Å) is normally incident on a slit. The central maximum fans out at
30° on both sides of the direction of the incident light. Calculate the slit width. For what width of the slit the
central maximum would spread out to 90° from the direction of the incident light?
Solution Given l = 5.0 ¥ 10–7 m and q = 30°.
l
b= (m = 1 for first minimum)
sin q
5.0 ¥ 10-7
= = 1.0 mm
sin 30∞
If q = 90°, b = ?
5.0 ¥ 10-7
b= = 0.5 mm
sin 90∞
E xamplE 18 A parallel beam of light (l = 5890 Å) is incident perpendicularly on a slit of width 0.1 mm.
Calculate angular width and linear width of central maximum formed on the screen 100 cm away.
Solution Given l = 5.89 ¥ 10–7 m, b = 1.0 ¥ 10–4 m and D = 1.0 m.
For angular width of central maximum
l
sin q = (m = 1 for first minimum)
b
5.89 ¥ 10-7
= = 5.89 ¥ 10–3
1.0 ¥ 10-4
q = sin–1 (0.00589)
= 0.3375°
Therefore, the total angular spread of central maximum is 2q, then
2q = 2 ¥ 0.3375°
2q = 0.675°
For linear width formula used is
x
sin q = [Please see Fig. 2.27]
D
x = D sin q = 1.0 ¥ 5.89 ¥ 10–3
= 5.89 ¥ 10–3 m
Total linear separation = 2x = 2 ¥ 5.89 ¥ 10–3 m
= 0.01178 m
= 1.178 cm
E xamplE 19 A single slit is illuminated by light composed of two wavelength l1 and l2. One observes
that due to Fraunhofer diffraction the first minimum obtained for l1 coincides with the second diffraction
minimum of l2. What is the relation between l1 and l2?
Solution Given m = 1 for l1 and m = 2 for l2.
For l1, m = 1, then
b sin q = l1 (i)
For l2, m = 2, then
b sin q = 2l2 (ii)
From Eqs. (i) and (ii)
l1 = 2l2
i.e., l1 is double of l2
E xamplE 20 Find the angular width of the central bright maximum in the Fraunhofer diffraction pattern of a
slit of width 12 ¥ 10–5 cm when the slit is illuminated by monochromatic light of wavelength 6000 Å.
l
or sin q = (m = 1)
b
6.0 ¥ 10-7
= = 0.5
1.2 ¥ 10-6
\ q = sin–1 (0.5) = 30°
or q = 30°
where q is half‑angular width.
Therefore, angular width of central maximum is
2q = 60°
E xamplE 21 Diffraction pattern of a single slit of width 0.5 cm is formed by a lens of focal length 40 cm.
Calculate the distance between the first dark and next bright fringe from the axis. Wavelength is 4890 Å.
Solution Given l = 4890 Å = 4.89 ¥ 10–7 m, f = 0.40 m and b = 0.005 m.
Formula used is b sin q = ml (i)
For first dark fringe, m = 1
l
sin q = (ii)
b
x
and sin q = [Please see Fig. 2.27] (iii)
f
By using Eqs. (iii) and (ii)
l x
=
b f
l f 4.89 ¥ 10-7 ¥ 0.40
x= =
b 5.0 ¥ 10-3
= 3.912 ¥ 10-5 m.
Diffraction 107
For first secondary maxima

(2n + 1)l
b sin q =
2
3l
sin q = (n = 1) (iv)
2b
x¢
and sin q = [from Fig. 2.27] (v)
f
x¢ 3l
=
f 2b
3l 3 ¥ 4.89 ¥ 10-7 ¥ 0.40
x¢ = f =
2b 2 ¥ 5.0 ¥ 10-3
or x¢ = 5.868 ¥ 10–5 m
Hence, the difference between 1st dark and next bright ring is
x2 – x1 = 5.868 ¥ 10–5 – 3.912 ¥ 10–5
Dx = 1.956 ¥ 10–5 m
E xamplE 22 A plane of wavelength 5893 Å passes through a slit, which is 0.5 mm wide and forms a
diffraction pattern on a screen placed on the focal plane of a lens of focal length 1.0 m. Calculate the separation
of the dark band on either side of the central maximum.
Formula used is b sin q = ml (m = 1)
l 5.893 ¥ 10-7 x
sin q = = = or x = 1.1786 ¥ 10–3 ¥ 1 = 1.1786 ¥ 10–3 m
b 5.0 ¥ 10-4 f
x = 1.1786 mm
Hence the separation of dark band on either side
= 2x = 2.357 mm.
E xamplE 23 Calculate the missing orders in a double slit Fraunhofer diffraction pattern, if the widths of slits
are 0.08 ¥ 10–3 m and they are 0.4 ¥ 10–3 m apart.
Solution Given b = 0.08 ¥ 10–3 and d = 0.4 ¥ 10–3 m.
The directions of interference maxima are given by
The directions of diffraction minima are given by
b sin q = ml (ii)
Dividing Eq. (i) by Eq. (ii), we get
b+d n (0.08 + 0.4) ¥ 10-3 n
= or =
b m 0.08 ¥ 10-3 m
6 n
or =
1 m
n = 6m
= 6, 12, 18, … etc. (m = 1, 2, 3, …)
th th th
Hence, 6 , 12 , 18 , … etc interference maxima will be missing in the diffraction pattern.
E xamplE 24 In a double‑slit Fraunhofer diffraction pattern the screen is 1.6 m away from the slits. The slit
width are 0.2 mm and they are 0.4 mm. Calculate the wavelength of light if the fringe width is 2.5 ¥ 10–3 m
and also deduce the missing orders.
Solution Given slit width b = 2.0 ¥ 10–4 m, 2d = 4.0 ¥ 10–4 m, b = 2.5 ¥ 10–3 m and D = 1.6 m.
b 2d 2.5 ¥ 10-3 ¥ 4.0 ¥ 10-4
Formula used is l= = = 6250 Å
D 1.6
The directions of interference maxima are
The directions of diffraction minima are
b sin q = ml (ii)
then, from Eqs (i) and (ii), we get
b+d n b+d
= or n = m
b m b
(2.0 + 4.0) ¥ 10-4
n= m = 3m
2.0 ¥ 10-4
= 3, 6, 9 … etc. (m = 1, 2, 3 …)
Hence 3rd, 6th, 9th, … interference maxima will be missing in the diffraction pattern.
E xamplE 25 A parallel beam of sodium light is normally incident on a plane transmission grating having
4250 lines per cm and a second order spectral line is observed at an angle of 30°. Calculate the wavelength
of light.
Solution Given N = 4250 lines per cm, q = 30° and n = 2.
Formula used is (b + d) sin q = nl
(b + d ) sin q
l=
n
1
Now, (b + d ) = cm
4250
1 sin 30∞
\ l= ¥
4250 2
= 5882 ¥ 10–8 cm
l = 5882 Å
E xamplE 26 A parallel beam of monochromatic light is allowed to incident normally on a plane transmission
grating having 5000 lines per cm and second order spectral line is found to be diffracted through 30°. Calculate
the wavelength of light.
Solution Given N = 5000 lines per cm, q = 30° and n = 2.
1 1
(b + d) = = = 2.0 ¥ 10–4 cm
N 5000
(b + d ) sin q
Formula used is (b = d) sin q = nl or l =
n
2.0 ¥ 10-4 ¥ sin 30∞
or l= = 5000 ¥ 10–8 cm
2
l = 5000 Å
Diffraction 109
E xamplE 27 In a grating spectrum, which spectral line in 4th order will overlap with 3rd order line of 5461 Å.
Solution Given n1 = 4, l2 = 5461 Å and n2 = 3, l1 = ?
As per question
(b + d) sin q = 4 l1 = 3 l2
3 3
or l1 = l2 = ¥ 5461 = 4096
4 4
or l1 = 4096 Å
E xamplE 28 In a plane transmission, grating the angle of diffraction for second order maxima for wavelength
5 ¥ 10–5 cm is 30°. Calculate the number of lines in one centimeter of the grating surface.
Solution Given l = 5 ¥ 10–5 cm, q = 30° and N = ?, n = 2.
(b + d) sin q = nl
nl 2 ¥ 5 ¥ 10-5 10 ¥ 10-5
(b + d ) = = =
sin q sin 30∞ 0.5
= 2.0 ¥ 10–4 cm.
The number of lines are
1 1 104
N= = -4
= per cm
(b + d ) 2.0 ¥ 10 2
= 5000 lines/cm
E xamplE 29 A plane grating has 15000 lines per inch. Find the angle of separation of the 5048 Å and 5016 Å
lines of helium in second order spectrum.
Solution Given l1 = 5048 ¥ 10–8 cm, l2 = 5016 ¥ 10–8 cm.
2.54
n = 2 and b + d = = 1.693 ¥ 10–4 cm
15000
For wavelength l1
(b + d) sin q1 = nl1
nl1 2 ¥ 5048 ¥ 10-8
sin q1 = =
b+d 1.693 ¥ 10-4
sin q1 = 0.5963
q1 = 36.60°
2 ¥ 5016 ¥ 10-8
Similarly, q 2 = sin -1
1.693 ¥ 10-4
= 36.34∞
Therefore, angle of separation Dq = q1 – q2
= 0.26°
E xamplE 30 A plane transmission grating having 6000 lines/cm is used to obtain a spectrum of light from
a sodium lamp in the second order. Calculate the angular separation between the two sodium lines whose
wavelengths are 5890 Å and 5896 Å.
1 1
Solution Given N = 6000 lines/cm, a + b + = cm, = 1.67 ¥ 10–4 cm.
N 6000
l1 = 5890 ¥ 10–8 cm, l2 = 5896 ¥ 10–8 cm, and n = 2
For l1
(b + d) sin q1 = nl1
nl1 2 ¥ 5890 ¥ 10-8

sin q1 = =
(b + d ) 1.67 ¥ 10-4
sin q1 = 0.7054
q1 = sin–1 (0.7054)
= 44.86°
2 ¥ 5896 ¥ 10-8
Similarly, q2 = sin -1
1.67 ¥ 10-4
= sin–1 (0.7061)
= 44.92
Thus, the angular separation (Dq) = q2 = q1
= 44.92 – 44.86
q2 – q1 = 0.06°
E xamplE 31 A diffraction grating used at normal incidence gives a line (5400 Å) in a certain order superposed
on the violet line (4050 Å) of the next higher order. How many lines per cm are there in the grating if angle
of diffraction is 30°?
Solution Given l1 = 5400 ¥ 10–8 cm of nth order,
l2 = 4050 ¥ 10–8 cm of (n + 1)th order
Formula used is (b + d) sin q = nl1 = (n + 1)l2 (i)
(b + d )sin q
n= (ii)
l1
(b + d )sin q
and n +1= (iii)
l2
Eliminating n by using Eqs (ii) and (iii), we get
Ê 1 1ˆ
ÁË l - l ˜¯ (b + d) sin q = 1
2 1
l1l2 1 5400 ¥ 10-8 ¥ 4050 ¥ 10-8
(b + d) = = = 32400 ¥ 10–8 cm
l1 - l2 sin q (5400 - 4050) ¥ 10-8 ¥ sin 30∞
1
or N= = 3086 lines/cm
b+d
E xamplE 32 A plane transmission grating produces an angular separation of 0.01 radian between two
wavelengths observed at an angle of 30°. Given mean value of the wavelength as 5000 Å. Calculate the
difference in two wavelengths if the spectrum is observed in the second order.
Diffraction 111
Solution Given dq = 0.01 radian, q= 30° and l = 5000 ¥ 10–8 cm.

Formula used is (b + d) sin q = nl (i)
By differentiating Eq. (i)
(b + d) cos q dq = ndl (ii)
Dividing Eq. (ii) by Eq. (i), we get
cos q dq d l
= or dl = l cot qdq
sin q l
dl = 5000 ¥ 10–8 cot 30° ¥ 0.01
dl = 86.6 Å
E xamplE 33 How many orders will be visible if the wavelength of the incident radiation is 5000 Å and the
number of lines on the grating is 2620 in one inch.
Solution Given N = 2620 lines per inch and l = 5000 ¥ 10–8 cm.
1
b+d= = 1 inch = 2.54 cm = 9.695 ¥ 10–4 cm
N
For maximum possible value sin q = 1, then
(b + d ) 9.695 ¥ 10-4
Order of spectrum (n) = =
l 5.0 ¥ 10-5 cm
= 19.38
= 19
That is 19th order will be visible.
E xamplE 34 What is the highest order spectrum which may be seen with monochromatic light of wavelength
5000Å by means of diffraction grating with 5000 lines/cm?
Solution Given N= 5000 lines per cm and l = 5000 ¥ 10–8.
1 1
b+d = = cm
N 5000
Formula used is (b+ d) sin q = nl (i)
For highest order spectrum to be visible the value of sin q must be 1. Then Eq. (i) becomes
b+d 1 1 100
n= = = =4
l 5000 5 ¥ 10-5 25
that is the highest order will be 4.
E xamplE 35 How many orders will be observed by a grating having 4000 lines per cm if it is illuminated by
visible light in the range 4000 Å to 7000 Å.
1
Solution Given (b + d) = cm = 2.5 ¥ 10–4 cm = 2.5 ¥ 10–6 m.
4000
For l1 = 4000 Å
(b + d )
n1 = (sin q = 1)
l1
2.5 ¥ 10-6
n1 = = 6.25
4.0 ¥ 10-7
n1 = 6.25
For l2 = 7000 Å
2.5 ¥ 10-6
n2 = = 3.57
7.0 ¥ 10-7
\ The order of the spectrum varies from 3 to 6 depending upon the wavelength of the visible range.
E xamplE 36 A diffraction grating having 4000 lines/cm is illuminated normally by light of wavelength 5000 Å.
Calculate its dispersive power in third order spectrum.
Solution Given l = 5.0 ¥ 10–5 cm, n = 3, N = 4000 lines/cm and
1
(b + d) = cm = 2.5 ¥ 10–4 cm.
4000
dq n
Formula used is = (i)
d l (b + d ) cos q
dq n
or = (ii)
d l (b + d ) 1 - sin 2 q
nl
As (b + d) sin q = nl or sin q = (iii)
(b + d )
By using Eqs (ii) and (iii), we have
dq n 3
= =
dl Ï nl ¸
2
Ê 3 ¥ 5 ¥ 10-5 ˆ
2
Ê 1 ˆ
(b + d ) 1 - Ì ˝ ÁË ˜¯ 1 - Á
Ó (b + d ) ˛ 4000 Ë 1 / 4000 ˜¯
12000
=
1 - 0.36
dq
or = 1.875 ¥ 104 rad/cm
dl
E xamplE 37 Calculate the minimum number of lines in grating which will first resolve the lines of
wavelengths 5890 Å and 5896 Å in the second order.
Solution Given l1 = 5.89 ¥ 10–5 cm, l2 = 5.896 ¥ 10–5 cm, n = 2.
\ Dl = 6 Å = 6 ¥ 10–8 cm
l
Resolving Power = = nN
dl
1 l 5.890 ¥ 10-5 5890
N= = = = 490.8
n Dl 2 ¥ 6 ¥ 10-8 12
E xamplE 38 For proper resolution 491 lines are required. Calculate the minimum number of lines in a
grating which will just resolve the sodium lines in the first order spectrum. The wavelengths are 5890 Å and
5896 Å.
Solution Given l1 = 5.89 ¥ 10–5 cm, l2 = 5.896 ¥ 10–5 cm and
Dl = 6.0 ¥ 10–8 cm, n = 1.
Diffraction 113
Formula used is l/Dl = n/N

l 5.89 ¥ 10-5
N= = = 981.66
nd l 1 ¥ 6 ¥ 10-8
i.e., for proper resolution the minimum number of lines in grating will be 982.
E xamplE 39 In a grating the sodium doublet (5890 Å, 5896 Å) is viewed in third order at 30° to the normal
and is resolved. Determine the grating spacing and the total width of the rulings.
Solution Given q = 30°, n = 3, lmean = 5.893 ¥ 10–5 cm.
nl 3 ¥ 5.893 ¥ 10-5 cm
(b + d ) = =
sin q sin 30∞
3 ¥ 5.893 ¥ 10-5
=
0.5
= 3.536 ¥ 10-4 cm
Q Resolving Power = n/N
l
\ = nN
dl
l 5.893 ¥ 10-5 cm
or N= = = N = 327
nd l 3 ¥ 6 ¥ 10-8 cm
Hence, the total width of the ruled surface is
N(b + d) = 327 ¥ 3.536 ¥ 10–4 = 0.1156 cm
E xamplE 40 Find the separation of two points on the moon that can be resolved by a 500 cm telescope. The
distance of the moon is 3.8 ¥ 105 km. The eye is most sensitive to light of wavelength 5500 Å.
Solution Given l = 5.5 ¥ 10–7 m, a = 5.0 m, R = 3.8 ¥ 108 m.
Limit of resolution of a telescope
1.22l 1.22 ¥ 5.5 ¥ 10-7
q= =
a 5.0
= 1.342 ¥ 10–7 rad
Let x be the distance between points, then
x
q=
R
x
1.342 ¥ 10–7=
3.8 ¥ 108
or x = 50.996 m
E xamplE 41 What will be the diameter of a telescope objective which is required to resolve two stars
separated by an angle of 10–3 degree? Assume l = 500 Å.
p
Solution Given l = 5.0 ¥ 10–7 m, q = 10–3 deg = ¥ 10–3 rad.
180
1.22l
Formula used is q=
a
1.22l 1.22 ¥ 5.0 ¥ 10-7

or a= =
q p
¥ 10-3
180
1.22 ¥ 1.8 ¥ 50 ¥ 10-3
= = 0.03495 m
3.14
E xamplE 42 Calculate the aperture of the objective of a telescope which may be used to resolve two stars
separated by 2.44 ¥ 10–6 radian for light of wavelength 6000 Å.
Solution Given l = 6.0 ¥ 10–7 m and q = 2.44 ¥ 10–6 rad.
1.22l 1.22l
Formula used is q= or a =
a q
1.22 ¥ 6.0 ¥ 10-7
or a= = 0.30 m
2.44 ¥ 10-6
or a = 0.30 m
Hence, aperture of the objective is 0.30 m.
E xamplE 43 Two pin holes 1.5 mm apart are placed in front of a source of light of wavelength = 5.5 ¥ 10–5
cm and seen through a telescope with its objective stopped down to a diameter of 0.4 cm. Find the maximum
distance from the telescope at which the pin holes can be resolved.
Solution Given l = 5.5 ¥ 10–5 m, a = 0.004 m and x = 1.5 ¥ 10–3 m.
1.22l
Formula used is q= (i)
a
x
and q= (ii)
R
x 1.22l
=
R a
xa 1.5 ¥ 10-3 ¥ 4.0 ¥ 10-3
or R= =
1.22l 1.22 ¥ 5.5 ¥ 10-7
= 8.9418 m
E xamplE 44 A microscope objective gathers light over a come of semi‑angle 30° and uses visible light
(l = 5500 Å). Estimate its resolving limit.
Solution Given l = 5.5 ¥ 10–7 m and q = 30°.
Formula used is 2d sin q = 1.22 l
1.22l
or d=
2 sin q
1.22 ¥ 5.5 ¥ 10-7
d= = 6.76 ¥ 10 –7 m
2 ¥ sin 30∞
Thus, the resolving limit of the microscope is 6.7 ¥ 10–7 m.
E xamplE 45 A microscope is used to resolve two self‑luminous objects separated by a distance of 4.0 ¥ 10–5
cm. If the wavelength of light is 5461 Å, compute the numerical aperture of the objective.
Diffraction 115
Solution Given l = 5.461 ¥ 10–7 m and d = 4.0 ¥ 10–7 m.

1.22l
Formula used is Numerical aperture (NA) =
2d
1.22 ¥ 5.461 ¥ 10-7
NA = = 0.833
2 ¥ 4.0 ¥ 10-7
Numerical aperture of the objective is 0.833.
E xamplE 46 A plane wave of light of wavelength 690 nm is incident on a vertical slit of width 10–4 m.
Sketch the intensity distribution on a screen 3 m from the slit placed parallel to the slit aperture. At what
distances from the central maximum do the first two zeroes occur?
Solution Given l = 6.90 ¥ 10–7 m, b = 10–4 m, D = 3.0 m.
x
By using the relation, b sin q = nl and sin q =
D
nl D
we get x=
b
for n=1
1 ¥ 6.90 ¥ 10-7 ¥ 3.0
x= = 2.07 cm
10-4
for n=2
2 ¥ 6.90 ¥ 10-7 ¥ 3.0
x= = 4.14 cm
10-4
E xamplE 47 A grating having 15000 lines per inch produces spectra of a mercury arc. The green line of the
mercury spectrum has a wavelength of 5461 Å. What is the angular separation between the first order and
second order green line?
Solution Given N = 15000 lines per inch, l = 5461 ¥ 10–8 cm.
By using the relation (b + d) sin q = nl
2.54
b+d= cm = 1.6933 ¥ 10–4 cm
15000
nl n ¥ 5461 ¥ 10-8
then sin q = =
b+d 1.6933 ¥ 10-4
for n = 1, q = q1
1 ¥ 5461 ¥ 10-8
sin q1 = = 0.322506
1.6933 ¥ 10-4
q1 = sin–1 (0.322506) = 18.81°
for n = 2, q = q2, then
nl 2 ¥ 5461 ¥ 10-8
sin q2 = =
(b + d ) 16933 ¥ 10-8
= 0.645013
q2 = sin–1 (0.645013) = 40.17°
Therefore, the angular separation of lines (Dq)
= q 2 – q1
= 40.17° – 18.81° = 21.36°
E xamplE 48 Light is incident normally on a grating 0.5 cm wide with 2500 lines. Find the angle of the two
sodium lines in the first order spectrum. Are the two lines resolved?
0.5 cm
Solution Given b + d = = 2 ¥ 10–4 cm, l1 = 5890 Å.
2500
and l2 = 5896 Å
q – q1 for D1, lines (n = 0)
Ï nl ¸
q1 = sin–1 Ì ˝
Ó (b + d ) ˛
È l1 ˘ È 5890 ¥ 10-8 ˘
q1 = sin–1 Í ˙ = sin -1 Í -8 ˙
= sin–1 (0.2945)
Î (b + d ) ˚ Î 20000 ¥ 10 ˚
q1 = 17.13°
q2 for D2 line (n = 1)
È l2 ˘ È 5896 ¥ 10-8 ˘
q2 = sin–1 Í ˙ = sin -1 Í -8 ˙
Î (b + d ) ˚ Î 20000 ¥ 10 ˚
= sin–1 [0.2948] = 17.15°
Dq = q2 – q1 = 0.02° = 0.02 ¥ 60¢ = 1.2¢ ¥ 60≤ = 72≤
Yes, these two lines will be just resolved.
Q.1 The diffraction phenomenon is

(a) bending of light around an obstacle (b) rectilinear propagation of light
(c) oscillation of light wave in one direction (d) none of them.
Q.2 The diameter of half‑period zones are proportional to
1
(a) 1/ n (b) (c) n (d) 1/n
n
Q.3 The ratio of areas of 1st to IInd half‑period zones is
(a) 1 : 2 (b) 1: 1 (c) 2: 1 (d) 4 : 1
Q.4 To find prominent diffraction, the size of the diffracting object should be
(a) greater than the wavelength of light used (b) of the order of wavelength of light
(c) less than the wavelength of light (d) none of these
Q.5 Product of focal lengths of a zone plate and a convex lens
(a) is independent of wavelength of light
(b) is directly proportional to wavelength of light
(c) is inversely proportional to wavelength of light
(d) none of the above
Q.6 Zone plate has some similarities as well as some differences with a
(a) plano‑convex lens (b) concave lens
(c) grating (d) convex lens
Diffraction 117
Q.7 In Fresnel diffraction

(a) source of light is kept at infinite distance from the aperture
(b) source of light is kept at finite distance from the aperture
(c) convex lens is used
(d) aperture width is selected so that it can act as a point source
Q.8 In Fraunhofer diffraction, the incident wavefront should be
(a) elliptical (b) plane (c) spherical (d) cylindrical
Q.9 In the diffraction pattern using circular aperture, when the screen is brought toward the aperture
(a) the intensity at the screen gradually increases
(b) the intensity at the screen gradually decreases
(c) the light is found to focus only to a fixed distance
(d) many points are observed where greater intensity is found
Q.10 Significant diffraction of x‑rays can be obtained
(a) by a single slit (b) by a double slit
(c) by a diffraction (d) by an atomic crystal
Q.11 Maximum number of orders found with a grating is
(a) independent of grating element (b) directly proportional to the grating element
(c) proportional to the wavelength (d) none of these
Q.12 The resolving power of a grating increases
(a) as the number of ruled lines per cm of a grating increases
(b) as the number of ruled lines per cm of a grating decreases
(c) as grating element decreases
(d) none of these
Q.13 Which of the following statement is true in the case of dispersive power of a grating?
(a) it is directly proportional to number of lines per cm of grating
(b) it is inversely proportional to cos q
(c) it is directly proportional to order of diffraction
(d) all of these
Q.14 In a plane diffraction grating, width of the principal maxima increases
(a) as number of slits decreases (b) as number of slits increases
(c) both (a) and (b) (d) none of these
Q.15 Which of the following statement is correct, if the width of the slit is reduced in single slit experiment?
(a) the intensity of fringes larged (b) the thickness of fringes reduced
(c) the thickness of fringes increased (d) none of these
Q.16 Which of the following colour (wavelength) will be diffracted more, when white light is used?
(a) red (b) green (c) violet (d) none of these
Q.17 The resolving limit of normal eye is
Ê 1ˆ Ê 1ˆ Ê 1ˆ
(a) Á ˜ ∞ (b) ÁË ˜¯ ∞ (c) ÁË ˜¯ ∞ (d) none of these
Ë 60 ¯ 30 45
Q.18 Which of the following relation is true for the limit of resolution of telescope?
1 a 1.22l
(a) (b) (c) (d) (a) and (b)
q 1.22l a
Q.1 What do you understand by diffraction of light?

Q.2 How will you distinguish between interference and diffraction?
Q.3 Discuss difference between Fraunhofer and Fresnel diffraction.
Q.4 Explain why Ist secondary maximum corresponding to p/2 is not obtained in single and double slit
diffraction pattern?
Q.5 What are Fresnel’s half‑period zones?
Q.6 What are the radii of zones of zone plate?
Q.7 In what respect a zone plate is similar or differ from a convex lens.
Q.8 What are the conditions of secondary maxima in diffraction due to single slit?
Q.9 What is a diffracting grating?
Q.10 On what factors does the dispersive power of a grating depend?
Q.11 What doe you mean by dispersive power?
Q.12 What is the ratio of intensity of secondary maximum and principal maximum?
Q.13 What is Rayleigh criterion of resolution?
Q.14 What do you mean by resolving power of a telescope?
Q.15 What do you mean by resolving power of a grating?
Q.16 What is the advantage of increasing the number of lines in a grating?
P RACtiCe P RoBleMs
general Questions
Q.1 How many types of diffractions are there? Distinguish between Fresnel and Fraunhofer type of
diffractions when the secondary wavelets are in the same phase at all points in the plane of the aperture.
Q.2 What is diffraction? Explain clearly the difference between interference and diffraction.
Q.3 What are Fresnel’s half‑period zones? Prove that the area of a half‑period zone on a plane wavefront is
independent of the order of the zone and that the amplitude due to a large wavefront at a point in front
of it is just half that due to the first half‑period zone acting alone. Hence give Fresnel’s explanations of
the rectilinear propagation of light.
Q.4 Explain the meaning of Fresnel’s half‑period zones. Why are they called so? What is the phase
difference between wavelets from successive half‑period zones? What are the factors on which the
amplitude of the light waves from a half‑period zone at the observation point depend?
Q.5 What is a zone plate and how is it made? Explain how a zone plate acts like a convergent lens having
multiple foci. Derive an expression for its focal length.
Q.6 How is zone plate constructed?
Q.7 Compare the performance of zone plate with that of a converging lens.
Q.8 Give the theory of a zone plate. Show that a zone plate has multi foci. Compare the zone plate with a
convex lens. What is meant by ‘phase reversal zone plate’?
Diffraction 119
Q.9 Give the mathematical treatment of Fresnel’s diffraction at a circular aperture.

Q.10 Describe the diffraction effects due to narrow circular aperture illuminated by monochromatic light
from a point‑source. Show that if the screen on which diffraction pattern is obtained be moving towards
the aperture, the intensity at the centre would be alternately maximum and minimum.
Q.11 Derive an expression for the intensity distribution due to Fraunhofer diffraction at a single slit and show
that the intensity of the first subsidiary maximum is about 4.5% of that of the principal maximum.
Q.12 Discuss the phenomena of Fraunhofer diffraction at a single slit and show that the relative intensities
of the successive maximum are nearly
1 : 4/9p2 : 4/25 p2 : 4/49 p2.
Q.13 Distinguish between interference and diffraction and explain the spectra formed by a grating and also
discuss the cause of decrease of intensity of principal maxima with the increase of order.
Q.14 If the width of the opaque space of grating is doubled the width of the transparent space, which order
of spectrum will be absent?
Q.15 Give the construction and theory of plane transmission grating and explain the formation of spectra by it.
Q.16 What is a grating? Explain the spectra, with theory, formed by a plane transmission diffraction grating.
Show that the intensity is not uniformly distributed over all the maxima.
Q.17 Define resolving power and dispersive power of a grating. Obtain expressions for these in the case of
plane transmission grating.
Q.18 Explain the phenomenon of Fraunhofer diffraction at a single slit.
Q.19 Explain how a chromatic beam gets dispersed while passing through a plane transmission diffraction
grating. Derive an expression for the dispersive power.
Q.20 Write short note on plane transmission grating.
Q.21 Deduce expressions for dispersive and resolving power of a diffraction grating. Can we increase
resolving power without affecting dispersive power? Discuss.
Q.22 Discuss two methods in detail for finding the wavelength of a given monochromatic light‑one using
interference of light and another based upon diffraction of light phenomenon.
Q.1 A screen is placed at a distance of 100 cm from a circular hole illuminated by parallel beam of light of
wavelength 6000 Å. Compute the radius of the fourth half‑period zone. [Ans: 1.6 ¥ 10–3 m]
Q.2 Find the first three focal lengths of a zone plate for which the radius of the first zone is 0.3 mm, for
light of wavelength 5000 Å. [Ans: 0.18 m, 0.06 m, 0.036 m]
Q.3 What is the radius of the first zone in a zone plate of focal length 20 cm for light of wavelength
5000 Å. [Ans: 3.16 ¥ 10–3 m]
Q.4 An object is placed at 20 cm from the zone plate and the brightest image is situated at 20 cm from the
zone plate with light of l = 4000 Å. Calculate the number of Fresnel’s zones in the radius of 1.0 cm of
the plate. [Ans. 2500]
Q.5 A zone plate is made such that the radii of the circles defining the zones are the same as the radii of
Newton’s rings formed between a plane surface and a surface whose radius of curvature is 150 cm.
Find the primary focal length of the zone plate. [Ans. 1.5 m]
Q.6 When a circular aperture of diameter 2.0 mm is illuminated by a plane wave of light, the most intense
point on the axis is at a distance of 200 cm from the aperture. Calculate the wavelength of light.
[Ans. 5000Å]
Q.7 A light of wavelength 6000 Å passes through a narrow aperture of radius 0.09 cm. At what distance
along the axis will first maximum intensity be observed? [Ans. 1.35 m]
Q.8 Calculate the half‑angular width of the central bright maximum in Fraunhofer diffraction at a slit
having width 1.23 ¥ 10–6 m illuminated by monochromatic light of wavelength 5896 Å[Ans. 28.61°]
Q.9 A parallel beam of light of wavelength 6000 Å falls normally on a slit. In the diffraction pattern the
first minimum lies at a distance 4.9 mm from the central maximum and the screen is placed 1.0 m from
the slit. Calculate the width of the slit. [Ans. 0.122 mm]
Q.10 A single slit of width 0.14 mm is illuminated normally by monochromatic light and diffraction bands
are observed on a screen 2.0 m away. If the centre of second dark band is 1.6 cm from the middle of
the central bright band, deduce the wavelength of light used. [Ans. 5600Å]
Q.11 A screen is placed 2.0 m away from a narrow slit which is illuminated with light of wavelength 6000Å.
If the first minimum lies 5 mm on either side of the central maximum. Calculate the slit width.
[Ans. 0.24 mm]
Q.12 A parallel beam of monochromatic light is normally incident on a plane transmission grating having
12000 lines per cm. The second order spectral line is observed at an angle 45°. Find the wavelength of
light used. [Ans. 2946 Å]
Q.13 A plane transmission grating having 5500 lines per cm is used to produce a spectrum of mercury light.
What will be the angular separation of the two yellow lines 4770 Å and 5790 Å in the second order.
[Ans. 10¢ of acc]
Q.14 A parallel beam of monochromatic light is allowed to be incident normally. On a plane grating having
1250 lines per cm and a second order spectrum line is observed to be deviated through 30°. Calculate
the wavelength of the spectral line. [Ans. 2 ¥ 10–4 cm]
Q.15 Find the angular separation between two sodium lines 5890 Å and 5896 Å in the second order spectrum
of a grating with 5000 lines/cm. The width of a grating is 0.5 cm. Can they be seen distinctly?
[Ans. 3 min, yes]
–5
Q.16 A plane transmission grating has 40,000 lines in all, with grating element 12.5 ¥ 10 cm. Calculate the
maximum resolving power for which it can be used in the range of wavelength 5000 Å.[Ans. 80,000]
Q.17 Calculate the aperture of the objective of a telescope which may be used to resolve stars separated by
4.88 ¥ 10–6 radian for light of wavelength 6000 Å. [Ans. 0.15 m]
Q.18 A telescope objective has a focal length of 3.0 m and a diameter of 0.01 m. Find the distance between
centres of the images of the two stars which are just resolved by it, assuming the wavelength of the
light 5000 Å. [Ans. 2.01 ¥ 10–4 m]
Q.19 Calculate the resolving power of a laboratory microscope if N.A. given on the objective is 0.12 and the
wavelength of light used is 6000 Å. [Ans. 4000]
–7
Q.20 A microscope is used to resolve two equally bright point objects separated by 5.55 ¥ 10 m. Calculate
the numerical aperture of the objective if light of wavelength 5460 Å is used. [Ans. 0.6]
Polarisation 3
Learning Objectives
After reading this chapter you will be able to
LO1 Explain transverse wave nature of LO5 Analyse optical activity and
polarisation phenomenon of specific rotation
LO2 Illustrate the difference between LO6 Demonstrate working of Half-shade
unpolarised and polarised light polarimeter, Biquartz polarimeter and
LO3 Enable to know the means of producing saccharimeter
plane-polarised light–transmission, LO7 Analyse photoelasticity
reflection, refraction and scattering
LO4 Discuss theory of production of plane,
circularly and elliptically polarised light
Introduction
You would have encountered problem in receiving signal when you watch TV or listen to your stereo
system. In order to overcome this problem, you adjust the position of an antenna attached to the stereo
system or align the TV antenna (receiving antenna) in the proper orientation. Did you ever think why this
is necessary and what physics is involved in doing so? Actually this is required as some types of antennas,
via the electrons, respond to the electric field of an electromagnetic wave (signal). If the orientation of the
receiving antenna matches with the orientation of the electric field of the wave, the electric field causes
the electrons to flow along the wires to generate a current. So the plane of the receiving antenna must be
horizontal if the electric field of the signal broadcast by the station vibrates in a horizontal plane. If the field
vibrates in a vertical plane, the orientation of the antenna should be changed to the vertical plane. So this is
clear that by doing the adjustments of position or orientation of the antenna, we increase the strength of
the signal, i.e., we improve the reception of the signal. The proper orientation of vibration of the electric
field is nothing but the polarisation of the wave.
A light wave is an electromagnetic wave whose electric field and magnetic field vectors vibrate
perpendicular to the direction of wave propagation. In order to completely identify the electromagnetic
wave, it is sufficient to specify the electric field since the magnetic field can be determined once the
electric field is known (discussed later in the chapter on Electromagnetic Wave Propagation). So a light wave
whose electric field vector, also called as light vector, is vibrating in more than one plane is referred to
as unpolarised light. The light emitted by the sun, by a lamp, or by a candle flame is unpolarised light. It
is possible to convert unpolarised light into polarised light in which the vibrations occur only in a single
plane. The process of converting unpolarised light into polarised light is known as polarisation. There are
a variety of methods of polarising light. Any interaction of light with matter whose optical properties are
asymmetrical along the directions transverse to the propagation vector provides a means of polarising
light. Only transverse waves can be polarised. The polarisation of longitudinal waves such as sound waves
is not possible as in these waves the vibrations occur only in the direction of wave propagation.
The phenomena of interference and diffraction discussed in the previous chapters show that the light
travels in the form of waves. However, these phenomena do not tell us about the nature of light waves,
i.e., whether the light waves are transverse or longitudinal or whether the vibrations are linear, circular or
elliptical. Such important investigations represent the subject of polarisation of light.
3.1 MechanicaL experiMent ShOwing pOLariSatiOn LO1

Of tranSverSe wave
In order to explain the polarisation of light, we consider a mechanical experiment by using a string whose two
end points are P and Q and which is kept fixed with point Q. The string passes through two parallel slits S1 and
S2 (Fig. 3.1). When the end point P is shaken up and down parallel to the slit S1, a transverse wave is generated
that travels towards the point Q and reaches there if the slit S2 is parallel to the slit S1 (Fig. 3.1a). Now we rotate
the slit S2 through 90° (Fig. 3.1b). Under this situation, the vibrations of the string are completely stopped by
this slit and the string does not vibrate in the region between S2 and Q. As is clear from this experiment, the
vibrations of this wave are confined to a plane parallel to the axis of the slit S1. Therefore, this wave is called
plane-polarised wave or linearly polarised wave. If the end P of the string is moved in a circular manner, the
P Q
S1 (a) S2
P Q
S1 S2
(b)
Figure 3.1
Polarisation 123
particles of the string will vibrate in the circle and the wave thus generated is called circularly polarised wave.
If the end P of the string is moved in an elliptic manner, the particles of the string will vibrate in the ellipse and
the wave thus generated is called elliptically polarised wave. Under these situations, only those vibrations can
pass beyond the slit S2 which are parallel to the axis of the slit S2. However, the passage of longitudinal waves
through the slit S2 will always be possible in its any position with respect to the slit S1.
3.2 Difference between UnpOLariSeD Light LO2

anD pOLariSeD Light
An unpolarised light is symmetrical about the direction of propagation as the light vector vibrates along all
possible directions perpendicular to the direction of propagation of light (Fig. 3.2a). In polarised light lack of
symmetry is found about the direction of propagation of the light. Since light is an electromagnetic wave, the
polarisation describes the orientation of the vibrations of the wave. The vibrations may be oriented in a single
direction in linearly polarised light (Fig. 3.2b, c). In Fig. 3.2b, the vibrations of light vector are in vertical
direction (in the plane of paper) whereas in Fig. 3.2c the vibrations are taking place in horizontal direction
(perpendicular to the plane of paper). However, in circularly polarised light and elliptically polarised light,
the direction of vibrations may rotate as the wave propagates. Circularly polarised light can rotate rightward
or leftward in the direction of wave propagation. Which of these two rotations is present in a wave is called
the wave’s chirality. Chirality is a Greek word which means handedness. This is a property of asymmetry
which is significantly useful in several branches of science.
(a) (b) (c)

Figure 3.2
3.2.1 plane of polarisation and plane of vibration

As discussed earlier, the electric field and magnetic field vectors of the light waves are perpendicular to each
other
 and to the direction of propagation of the wave. As per convention, the direction of electric field vector
E is taken
 as the direction of polarisation of the wave. In view of this, the plane determined by the electric
vector E and the direction of wave propagation (vector k ), is called the plane of polarisation, particularly
for the plane (linearly) polarised light. This term is sometimes applied to the plane at right angles to this,
i.e., to the plane
 containing the magnetic field vector. Then the plane containing the electric vector E and
wave vector k is named as the plane of vibration. However, in a plane electromagnetic wave with elliptic
polarisation and circular polarisation,
 the electric field vector E has components
 (out of phase) in both the
directions perpendicular to k . So the plane that contains only vibrations of E field can be called as the plane
of polarisation. Hence, the plane of vibration and the plane of polarisation will be the same in this case.
3.3 Means of Production of Plane–Polarised light LO3

The different methods commonly used for the production of plane-polarised light may be classified under
different categories: (a) polarisation by transmission, (b) polarisation by reflection, (c) polarisation by
refraction, and (d) polarisation by scattering.
3.3.1 Polarisation by transmission
It is possible to get the unpolarised light converted into a polarised light along with its vibrations in a single
plane if the unpolarised light is passed through a Polaroid filter.
3.3.1.1 Polaroid Filter
The most general method of polarisation makes use of a Polaroid filter. The Polaroid filter has long chain
molecules that are aligned in the same direction within the filter. The alignment of these molecules constitutes
a polarisation axis that extends across the length of the filter. This axis allows electromagnetic waves to pass
through whose vibrations are parallel to the axis. Thus any vibrations perpendicular to the polarisation axis
are stopped by the filter. When an unpolarised light is passed through a Polaroid filter, it emerges with its
vibrations in a single plane along with one half of its intensity. This way the emerging light is the polarised
light. The relationship between the alignment of long chain molecules and the polarisation axis in a Polaroid
filter is just opposite, i.e., a Polaroid filter with its long chain molecules aligned vertically will have a
horizontally aligned polarisation axis. This type of a filter will stop all the vertical vibrations and allow only
the horizontal vibrations to pass through. However, a Polaroid filter with its long chain molecules aligned
horizontally will have a polarisation axis aligned vertically. So this filter will stop all the horizontal vibrations
and allow only the vertical vibrations to pass through. It is clear that the two Polaroid filters oriented with
their polarisation axes perpendicular to each other will stop all the light.
Unpolarised Polaroid Polarised

Light Filter Light
Figure 3.3
3.3.1.2 Malus’ Law

Based on Malus’ cos2 q intensity law the intensity of transmitted light is described when there is a relative
orientation between the directions of the polarisation of the incoming light and the polarisation axis of the
filter. Therefore, the incoming light should be polarised. For this we first polarise the unpolarised light by
passing it through a Polaroid filter, called the polariser. Then the polarised light is passed through a second
Polaroid filter, called the analyser. The transmission axis of the analyser makes an angle q with the plane of
polarisation of the incident light (Fig. 3.4). This way the emerging light is polarised in the same direction as
the transmission axis of the analyser and its intensity varies as the cosine square of the angle between the
planes of the transmission of analyser and the polariser. This is called Malus’ law.
In order to get the relation for resultant intensity, we consider A to be the amplitude of the plane-polarised
light and q as the angle between the planes of the polariser and the analyser (Fig. 3.5). The amplitude A can
be resolved into two components A cos q and A sin q. Only the component A cos q is transmitted through the
analyser. Then a measure of the intensity of the transmitted light will be
Polarisation 125
Angle q
Unpolarised Polariser Analyser

Light
Figure 3.4
I = (A cos q)2
= A2 cos2q = I0cos2q (I0 = A2)
I = I0 cos2q
where I0 is the maximum intensity. This equation is known as
Plane of Polariser
Malus’ law, which gives the intensity of transmitted light. It is
clear from this relation that the whole intensity is passed when
the planes of the polariser and the analyser are parallel (q = 0°). Plane of
However, the incident light is completely blocked when these axes Analyser
are perpendicular to each other. A
q
os
3.3.2 Polarisation by reflection
Ac
We can obtain partially or sometimes fully polarised light when As q
the light is reflected by the surface of an electrical insulator. In this in
q
case, the degree of polarisation depends on the angle of incidence
of the light and the refractive index of the reflecting material. In
some cases, the reflected light can be completely polarised parallel Figure 3.5
to the reflecting surface and perpendicular to the direction of the
light propagation.
3.3.2.1 Brewster’s Law
In 1808, Malus discovered a simple method for the polarisation of light by reflection. He found that when an
ordinary light is reflected from the surface of a glass plate the reflected and refracted light beams are partially
plane-polarised (Fig. 3.6a). This depends on the angle of incidence and at a particular angle of incidence
(57.5° for a glass surface) the reflected light is completely plane-polarised while the transmitted light is
partially polarised. This angle of incidence is known as polarising angle. Brewster also performed series of
experiments in 1911 for studying the polarisation of light by reflection at the surfaces of different media. He
also found that at a particular angle of incidence, the reflected light is completely polarised. The reflected
light is the component of incident light polarised normal to the plane of incidence and therefore parallel to

the surface (Fig. 3.6b) in view of the plane of incidence as the plane having vector k and unit vector normal
to the surface. The incident angle at which the reflected light is completely polarised is known as Brewster’s
angle or angle of polarisation (ip).
Ordinary Light Partially

(Unpolarised) N Polarised A N
B
ip ip
i i
O O
r r
m m
C
N´ N´
(a) (b)
Figure 3.6
According to Brewster, the refractive index m of the medium (Fig. 3.6b) is given by
sin i p
m = tan i p = (i)
cos i p
The above relation which says that the tangent of the angle of polarisation is numerically equal to the refractive
index of the medium is called Brewster’s law. If the light is propagating in a medium with refractive index m1
and is being partially reflected at the boundary with a medium of refractive index m2, Brewster’s law however
takes the following form
m2
tan i p =
m1
The above polarising angle ip is sometimes referred to as the Brewster angle of the material. As per Snell’s
law, for a glass surface in Fig. 3.6b we can write
sin i p
m= (ii)
sin r
From Eqs. (i) and (ii), we get
sin i p sin i p
= (iii)
cos i p sin r
From Eq. (iii), we have

Êp ˆ
cos i p = sin r = cos Á - r ˜
Ë2 ¯
p
ip = -r
2
p
ip + r =
2
Polarisation 127
As we know, –NON ¢ = p, then we have,

–NOB + –BOC + –CON ¢ = p
ip + –BOC + r = p
\ –BOC = p – (r + ip)
p
=p–
2
p
–BOC =
2
i.e., at the polarising angle, the reflected and refracted rays are at right angle to each other.
3.3.2.2 Biot’s Polariscope
Biot’s polariscope is used for producing and detecting the plane-polarised light by reflection and is shown in
Fig. 3.7. In the figure A and B are two parallel glass plates and their back surfaces are blackened to prevent
the transmitted light (Fig. 3.7a). Let a ray of unpolarised light PQ be incident at the polarising angle on the
polished face of glass plate A and it is reflected along QR. This ray QR is incident on the other glass plate B at
the polarising angle and hence gets reflected along RS. It is seen that the ray RS is completely plane-polarised
and its intensity is found to be maximum. The two glass plates A and B are known as the polariser and the
analyser, respectively. When the glass plate B is gradually rotated about QR keeping the angle of incidence
constant, the intensity of the reflected ray RS decreases and it becomes zero for 90° rotation of glass plate B,
as shown in Fig. 3.7(b). If the glass plate B is further rotated, the intensity of the ray RS becomes maximum
at 180°, minimum at 270° and again becomes maximum at 360°. By using this simple experiment, we can
obtain and detect plane-polarised light.
B
R
B
R
ip S
N N
P P
ip ip
Q Q
A A
(a) (b)
Figure 3.7
3.3.3 Polarisation by refraction
Polarisation can also take place by the refraction of light, which occurs when a beam of light passes from
one material into another material. Under this situation, the path of the light beam changes its direction at the
surface of the two materials and then the refracted beam acquires some degree of polarisation. Mostly, the
polarisation occurs in a plane perpendicular to the surface.
The light split into two beams upon entering the crystal and both the refracted light beams are polarised – one
in a direction parallel to the surface and the other in a direction perpendicular to the surface. Since these two
refracted rays are polarised with a perpendicular orientation, a polarising filter can be used to completely stop
one of the images.
3.3.3.1 Optic Axis
You would have learnt in the experiments using lenses that there exists a line which passes through the centre
of curvature of a lens surface such that the light rays are neither reflected nor refracted. This is called the optic
axis or the principal axis. Similar situation arises for a particular A
type of crystal such as calcite crystal or tourmaline crystal (Fig.
3.8). For these substances there exists a specific direction within
the crystal known as the optic axis or the principal axis, which is 102°
102°
determined by the atomic configuration of the crystal. The optic
axis of a calcite crystal is shown in Fig. 3.8 by the dotted line AB. 102°
Any ray of ordinary unpolarised light incident along the optic
axis or parallel to this axis does not split up into two rays. The
light ray gets split into two rays called as ordinary ray (O-ray) 78° 78°
and extraordinary ray (E-ray) only when it makes an angle with
the optic axis. It is observed that the ordinary and extraordinary
102°
rays propagate at the same speed along the optic axis. This is
true for any direction which is parallel to the optic axis. The
crystal in which only one such axis (direction) exists is called
uniaxial crystal. The examples of uniaxial crystals are calcite,
tourmaline and quartz. The crystal in which two directions exist
along which the speeds of O-ray and E-ray are the same is called B
biaxial crystal. The examples of the biaxial crystals are topaz AB: Optic Axis
and aragonite.
Figure 3.8
3.3.3.2 Principal Section of a Crystal
The plane containing the optic axis and the perpendicular to the pair of opposite faces of the crystal is
known as principal section for that pair of faces of the crystal. Since the crystal has six faces, for each pair of
opposite faces of the crystal, there are three principal sections.
3.3.3.3 Geometry of Calcite Crystal
The calcite or calcspar is the commonest crystalline form of calcium carbonate (CaCO3). It is also known as
Iceland spar. It is a colourless crystal which is transparent for visible and ultraviolet light. It occurs in nature
in different variety of crystal form, for example in the rhombohedral class of the hexagonal system. It breaks
readily into simple cleavage rhombohedrons, whose shape is shown in Fig. 3.8. It can be seen from the figure
that its each face is a parallelogram with angles as 78° and 102°. An interesting feature of calcite is that each
crystal can be made to slice or break along cleavage planes into two or more smaller crystals with faces that
are parallelograms with angles 71° and 109° (Fig. 3.9).
Polarisation 129
O
pt A
ic 109° 71°
Ax
is
i Q R
E-ray
ay re
nt R
ide
P Inc r0 109° S
O-ray
71°
B
Figure 3.9
3.3.3.4 Double Refraction

Double refraction is an optical property in which a single ray of unpolarised light when passes through a
uniaxial crystal gets split into two refracted rays, each propagating in a different direction. These two rays are
named as O-ray and E-ray. Therefore, the phenomenon of splitting of light into two rays is known as double
refraction.
If a narrow beam of unpolarised light PQ is incident on the face AB of the crystal at an angle i, it splits up into
O-ray and E-ray inside the crystal. O-ray makes an angle of refraction ro whereas the E-ray makes an angle
re with the optic axis inside the crystal (Fig. 3.9). Both the rays are plane-polarised and their vibrations are at
right angles to each other and to the direction of propagation. The refractive index for O-ray is observed to
be constant in all the directions whereas the refractive index for E-ray varies according to the direction taken.
This is because it has components that are both parallel and perpendicular to the optic axis of the crystal.
Since the speed of light wave in a medium is the ratio of its speed in vacuum and the index of refraction for
that wavelength, an E-ray can move either faster or slower than an O-ray. However, the velocity of O-ray
inside the crystal is same in all the directions while the E-ray travels in the crystal with different velocities in
different directions. For example, in a negative crystal such as calcite, where the refractive index for E-ray is
smaller than that for O-ray (me < m0; m0 = 1.6584, me = 1.4864), the velocity of E-ray is larger than the velocity
of O-ray in the crystal. On the other hand, in a positive crystal such as quartz, where the refractive index for
E-ray is greater than that for O-ray (me > m0), the velocity of E-ray is smaller than the velocity of O-ray in the
crystal.
In view of Fig. 3.9, the refractive indices for O-ray and E-ray can be written as
sin i sin i
m0 = and me =
sin ro sin re
Since the velocity of O-ray and hence the refractive index m0 inside the crystal is same in all the directions,
this ray obeys Snell’s Law. However, the E-ray does not obey Snell’s Law as it travels in the crystal with
different velocities in different directions, leading to different me in different directions.
3.3.3.5 Polarisation by Double Refraction
The polarisation of light by double refraction in calcite is demonstrated in Fig. 3.10 where AB and CD are the
principal sections of the two crystals. Here we rotate the second crystal and observe the following phenomena
related to the O-ray and E-ray separated by the crystals.
(a) In the case of parallel principal sections of the two crystals, two images O and E are seen in
Fig. 3.10a. The O-ray from both the crystals passes undeviated and emerges as O1-ray. However,
the E-ray passes the second crystal along a path parallel to its path inside the first crystal and finally
emerges as E1-ray. This happens when the thickness of both the crystals is the same. Hence, the
images O1 and E1 are separated by a distance equal to sum of the two displacements found in each
crystal, if used separately.
(b) If the second crystal is rotated about the incident light taking it as the axis and keeping the first
crystal fixed, the O-ray and the E-ray split separately into two rays. So the two new images O2 and
E2 are observed along with O1 and E1. If we further rotate the crystal, the images O1 and O2 remain
fixed whereas E1 and E2 rotate around O1 and O2, respectively. Under this situation, the intensity of
O1 and E1 decreases. When the principal section of the second crystal makes an angle of 45° with
the principal section of the first crystal, the four images of equal intensities are seen. This is shown
in Fig. 3.10b.
A C
A C
E E1
E1
O1
O1
O
B D
B D
(a)
A
A
C
O2 E1 O2 E1 O2
E2 D C O1
O1 E2 E2
D
B
B
(b) (c) (d)
A D
A D
O1 E1
E1
O1
B C
B C
q = 180 q = 180
(e)
Figure 3.10
On continuing the rotation, the intensities of images O1 and E1 get decreased and the intensities of
O2 and E2 get increased. At 90° rotation, the images O1 and E1 finally disappear and the new images
O2 and E2 acquire maximum intensities Fig. 3.10c.
(c) For the further rotation of the second crystal, the images O1 and E1 again appear and the intensities
of these images increase. Then the intensities of the images O2 and E2 decrease. At 135° angle of
rotation, the intensities of four images become equal, as shown in Fig. 3.10d.
Polarisation 131
(d) At q = 180°, the principal sections of both the crystals are again parallel. However, their optic axes
are oriented in the opposite directions (Fig. 3.10e). In this situation, the images O2 and E2 disappear
and the images O1 and E1 superimpose with each other to form a single image that emerges from the
second crystal.
Based on the above observations, this experiment demonstrates the polarisation of light. The first crystal
produces plane-polarised vibrations whereas the second crystal analyses these vibrations.
Explanation Let the principal sections AB and CD of the first A
and second crystals, respectively, be inclined at an angle q
E
(Fig. 3.11). A ray of ordinary unpolarised light splits into two C
plane-polarised rays after emerging from the first crystal. When
the O-ray vibrates perpendicular to the principal section AB, then E1
1
E
the E-ray vibrations are along the principle section. Let a be the
q=
amplitude of each ray, represented by NO and NE, respectively.
os
a
ac
On entering the second crystal, each of the O- and E-rays is split E2
into two components. The O-ray is split into two components as
2
E
q
q=
O1 = a cos q and E2 = a sin q whereas E-ray is split into two O
=a
in
2
sin a
as
components as E1 = a cos q and O2 = a sin q. Thus, a measure of q O
N
the intensities of O1 and E1 is (a cos q)2 while that of O2 and E2 is q
ac
(a sin q)2. Based on these expressions for the intensities, we can os
q=
now discuss the different cases: O
1
2 2
Case-1: At q = 0° and q = 180°, we get cos q = 1 and sin q = 0. It D B
O1
means the intensities of O1 and E1 rays are maximum, while that
of O2 and E2 is zero. Figure 3.11
2 2
Case-2: If q = 45° and q = 135°, we get cos q = ½ and sin q = ½. It means the intensities of O1, E1, O2 and
E2 have the same values. Therefore, all the four images are equally bright.
Case-3: If q = 90°, we get cos2q = 0 and sin2q = 1. It means O1 and E1 vanish and O2 and E2 are the brightest.
It is clear from the above cases that the sum of the intensities of the two components is a2cos2q + a2sin2q =
a2, which is just equal to the intensity of the incident light beam.
3.3.3.6 Huygens’ Theory of Double Refraction
Phenomenon of double refraction was explained by Huygens’ for which he extended his principle of secondary
wavelets and made the following assumptions.
(i) When a light wave strikes the surface of a doubly refracting crystal, each point of the crystal
becomes the origin of two secondary wavelets, named as ordinary ray and extraordinary ray. These
two wavelets spread out into the crystal.
(ii) The wavefront corresponding to ordinary ray is spherical as the velocity of ordinary ray remains the
same in all the directions (Fig. 3.12a).
(iii) The wavefront corresponding to extraordinary ray is an ellipsoid of revolution with the optic axis
as its axis of revolution (Fig. 3.12b). This is due to the fact that the velocity of E-ray is different in
different directions in the crystal.
(iv) The two wavefronts corresponding to O-ray and E-ray touch each other along the optic axis since
both the rays travel with the same velocity along the direction of optic axis.
(v) For negative uniaxial crystals (like calcite) in which the velocity of O-ray is less than the velocity
of E-ray, sphere lies inside the ellipsoid (Fig. 3.12c). However, for positive uniaxial crystals (like
quartz) the ellipsoid lies inside the sphere (Fig. 3.12d) since in this case the velocity of O-ray is
greater than the velocity of E-ray.
Optic Axis
(a) Spherical (b) Ellipsoid of (c) Negative (d) Positive

Wavefront Revolution Crystal Crystal
Figure 3.12
3.3.3.7 Nicol Prism

Nicol prism is an optical device which is used for producing and analysing plane-polarised light in practice.
It is constructed from the calcite crystal PQRS having length three times of its width. Its end faces PQ and RS
are cut such that the angles in the principal section become 68° and 112° in place of 71° and 109° (Fig. 3.13).
The crystal is then cut diagonally into two parts. The surfaces of these parts are grinded to make optically flat
and then these are polished. Thus polished surfaces are connected together with a special cement known as
Canada Balsam.
P¢
P R
E-ray
71° 112°
68° 109°
Q S
S¢
O-ray
Figure 3.13
3.3.3.8 Quarter-Wave Plate

It is a plate of doubly refracting uniaxial crystal like quartz or calcite, whose refracting faces are cut parallel
to the optic axis and its thickness is such that it introduces a phase change of p/2, i.e., a path change of l/4
between the ordinary and extraordinary light waves. Let m0 and me be the refractive indices for the ordinary
and extraordinary light waves, respectively, and t the thickness of the plate. The path difference between the
ordinary and extraordinary waves is given by
Polarisation 133
l
(m0 – me)t = for negative crystal
4
l
(me – m0)t = for positive crystal
4
l
t= for negative crystal
4( m0 - me )
Quarter-wave plate is used to produce circularly and elliptically polarised light.
3.3.3.9 Half-Wave Plate
It is a plate of doubly refracting uniaxial crystal like quartz or calcite, whose refracting faces are cut parallel
to the optic axis and its thickness is such that it introduces a phase change of p, i.e., a path change of l/2
between the ordinary and extraordinary light waves. For the refractive indices m0 and me for ordinary and
extraordinary light waves, the path difference is written as,
For Half-Wave Plate
l
(m0 – me)t = for negative crystal
2
l
(me – m0)t = for positive crystal
2
l
t= for negative crystal
2( m0 - me )
Half-wave plate is used to produce plane-polarised light. Quarter-wave plate and half-wave plate are known
as phase retarding plates. The phase retardation can be calculated by using the following relation.
2p
d= ¥ Dx where Dx is path difference.
l
When a beam of unpolarised light is incident on the face P ¢Q, it gets split into two refracted rays, named
O-ray and E-ray. These two rays are plane-polarised rays, whose vibrations are at right angles to each other.
The refractive index of Canada balsam cement being 1.55 lies between those of ordinary and extraordinary
rays. This is because the refractive indices of ordinary and extraordinary rays for calcite crystal are 1.6584
and 1.4864, respectively.
It is clear from the above discussion that Canada Balsam layer acts as an optically rarer medium for the
ordinary ray and it acts as an optically denser medium for the extraordinary ray. When ordinary ray of light
travels in the calcite crystal and enters the Canada balsam cement layer, it passes from denser to rarer medium.
Moreover, the angle of incidence is greater than the critical angle, the incident ray is totally internally reflected
from the crystal and only extraordinary ray is transmitted through the prism. Therefore, fully plane-polarised
wave is generated with the help of Nicol prism.
Nicol Prism as a Polariser and an Analyser: In order to produce and analyse the plane-polarised light,
we arrange two nicol prism as per Fig. 3.14. When a beam of unpolarised light is incident on the nicol
prism, emergent beam from the prism is obtained as plane-polarised, and which has vibrations parallel to the
principal section. This prism is therefore known as polariser. If this polarised beam falls on another parallel
nicol prism P2, whose principal section is parallel to that of P1, then the incident beam will behave as E-ray
P1 P2
E-ray
O-ray
(a)
P1 P2
O-ray
O-ray
(b)
P1 P2
E-ray
O-ray
(c)
Figure 3.14
inside the nicol prism P2 and gets completely transmitted through it (Fig. 3.14a). This way the intensity of
emergent light will be maximum.
Now the nicol prism P2 is rotated about its axis, then we note that the intensity of emerging light decreases
and becomes zero at 90° rotation of the second prism (Fig. 3.14b). In this position, the vibrations of E-ray
become perpendicular to the principal section of the analyser (nicol prism P2). Hence, this ray behaves as
O-ray for prism P2 and it is totally internally reflected by Canada balsam layer. This fact can be used for
detecting the plane-polarised light and the nicol prism P2 acts as an analyser.
If the nicol prism P2 is further rotated about its axis, the intensity of the light emerging from it increases and
becomes maximum for the position when principal section of P2 is again parallel to that of P1 (Fig. 3.14c).
Hence, the nicol prisms P1 and P2 acts as polariser and analyser, respectively.
3.3.4 Polarisation by scattering
The scattering of light by the air molecules produces linearly polarised light in the plane perpendicular to the
incident light. The scatterers can be imagined as tiny antennae which emit radiations (light) perpendicular to
their axis of vibrations. If the charges in a molecule are vibrating along the x-axis, the radiation or light is not
obtained along the x-axis rather the scattered light is found to be linearly polarised at 90° away from the beam
direction. This leads the light to be partially polarised that undergoes Rayleigh scattering from the blue sky.
Polarisation 135
3.4 theOry Of prODUctiOn Of pLane, circULarLy anD LO4

eLLipticaLLy pOLariSeD Light
As discussed earlier, the light that has unidirectional vibrations is known as plane-polarised light or linearly
polarised light. When two plane-polarised light waves are allowed to superimpose, and the resultant electric
vector rotates in such a way that its tip traces a circle, the resultant light is known as circularly polarised light.
However, if the magnitude of the resultant electric vector varies periodically during its rotation and its tip
traces an ellipse, then the resultant light is called elliptically polarised light.
Let us consider a calcite crystal such that its refracting faces are cut parallel to its optic axis. Further we
consider that the linear vibrations (amplitude A) in the incident light are along the direction PQ that makes an
angle q with the optic axis. (Fig. 3.15a). Under this situation, the incident plane-polarised light of amplitude A
spits into two components A cos q (E-ray) and A sin q (O-ray). These components constitute E-ray and O-ray
in view of their vibrations parallel to the optic axis and perpendicular to it, respectively. As per Huygens’
theory, the E-ray and O-ray travel in the same direction (Fig. 3.15b) with different velocities. Since calcite
is a negative crystal, the velocity of E-ray will be greater than that of O-ray. Hence, a phase difference is
introduced between them after traveling through the plate.
y
Optic Axis O-Ray Q
A P
Optic Axis
q
P x
E-Ray
O-Ray
E-Ray
(b)
(a)
Figure 3.15
In view of the incident light wave as A sin wt, we can represent the E-ray along the optic axis as
x = A cos q sin(wt + f) (i)
Similarly, the O-ray along y-axis will be
y = A sin q sin wt (ii)
Now assuming A cos q = a and A sin q = b, we get
x = a sin(wt + f) (iii)
y = b sin wt (iv)
From Eq. (iv), we have
y
sin wt = (v)
b
y2
and cos wt = 1- (vi)
b2
Now from Eq. (iii), we get
x
= sin wt cos f + cos wt sin f (vii)
a
Putting the values of sin wt and cos wt from Eqs. (v) and (vi) in the above equation, we have
x y y2
= cos f + 1 - 2 sin f
a b b
x y y2
or - cos f = 1 - 2 sin f (viii)
a b b
On squaring both sides, Eq. (viii) we get
˘ È y2 ˘
2
Èx y
-
ÍÎ a b cos f ˙˚ = ÍÎ1 - b 2 ˙˚ sin f
x2 y 2 2 2 xy È y2 ˘ 2
+ cos f - cos f = Í1 - ˙ sin f
a 2 b2 ab Î b2 ˚
x 2 y 2 2 xy
or + - cos f = sin 2 f (ix)
a 2 b2 ab
This is the general equation of an ellipse.
Special Cases: Since the phase difference f between the ordinary and extraordinary rays depend upon the
thickness of the plate, we will discuss below the different cases on the basis of this thickness t.
Case-I: If the thickness of the plate is such that it introduces a phase difference of f = 0, 2p, 4p, … between
O-ray and E-ray, then sin f = 0 and cos f = 1. Therefore, Eq. (ix) becomes,
x 2 y 2 2 xy
+ - =0
a 2 b2 ab
2
Èx y˘
or ÍÎ a - b ˙˚ = 0
b
or y= x (x)
a
This is the equation of straight line having the slope ÊÁ ˆ˜ and passing through the origin Fig. 3.16(a). This
b
Ë a¯
concludes that the light emerging through the plate is plane-polarised.
Case-II: If the thickness of the plate is such that f = p, 3p, 5p, ..., then sin f = 0 and cos f = –1. Therefore,
Eq. (ix) attains the form
Polarisation 137
y y
b b
O q q
a x a x
O q
(a) (b)
y
a
b
a x
a
(c)
(d)
Figure 3.16
x 2 y 2 2 xy
+ + =0
a 2 b2 ab
2
Èx y˘
or ÍÎ a + b ˙˚ = 0
-b
y= x (xi)
a
Ê bˆ
This is again an equation of straight line having the slope ÁË - ˜¯ (Fig. 3.16b). So we will have again the
a
emergent light as plane-polarised light.
p 3p 5p
Case-III: If the thickness of the plate is such that f = , , , then sin f = 1, cos f = 0 Eq. (ix) attains
2 2 2
the form
x2 y 2 (xii)
+ =1
a 2 b2
This is the equation of an ellipse with its axis along x and y directions (Fig. 3.16c). Therefore, the emergent
light will be elliptically polarised light.
Case-IV: If a = b and f satisfies the condition of Case-III
x2 + y2 = a2
This is the equation of a circle of radius a. Thus, the emergent light will be circularly polarised light if the
plate introduces a phase change of
p 3p 5p
, , etc.
2 2 2
From the above discussion it is clear that the plane and circularly polarised lights are the special cases of an
elliptically polarised light which is obtained by the superposition of two plane-polarised lights.
leavo-rotatory substance
dextro-rotatory substance
3.5 OpticaL activity LO5

When a beam of plane-polarised light is passed through certain substances, then the plane of polarisation of
the beam is rotated about the direction of propagation of the beam through a certain angle. This phenomenon
of rotating the plane of polarisation by certain substances is known as optical activity and the material is
known as optically active substance. The examples of optically active substances are sugar crystals, sugar
solution, sodium chlorate etc. There are two types of optically active substances; one rotates the plane of
polarisation to the right and the other rotates this plane to the left. The substances that rotate the plane of
polarisation to the right are known as dextro-rotatory or right handed, i.e., plane of polarisation is rotated in
clock wise direction from the point of view of the observer. The substances that rotate the plane of polarisation
to the left are known as leavo-rotatory or left handed, i.e., plane of polarisation is rotated in anticlockwise
direction from the point of view of the observer.
In 1815, Biot gave the following laws about the optical activity.
(i) The angle of rotation of a light beam of a given wavelength is directly proportional to the length of
optical active substance traversed.
(ii) For a given path length, the angle of rotation is proportional to the concentration of the solution or
vapour.
(iii) The rotation produced by optically active substances is equal to the algebraic sum of rotations made
by the substances individually. The anticlockwise and clockwise rotations are taken with opposite
signs.
(iv) The angle of rotation is inversly proportional to the square of wavelength. In case of quartz, the
angle of rotation is given by
Angle of rotation = A + B/l2. where A and B are constants.
3.6 Specific rOtatiOn LO5

It is an intrinsic property of a pure material at a given wavelength and temperature. The plane of the linearly
polarised light is rotated when it is passed through liquids containg an optically active substance, for example
sugar solution, sodium chlorate, cinnabar, camphor in alcohol, etc.). The specific rotation S is defined as the
observed angle of optical rotation q when plane-polarised light is passed through a sample with a path length
of 1 decimetre and a sample concentration of 1 gram per 1 millilitre. Therefore
q q qV
S= = =
l ¥ c l ¥ ( m /V ) l ¥ m
In this equation, l is the path length in decimeters, and c is the concentration of the liquid in g/ml, for a sample
at a temperature T(given in degrees Celsius) and wavelength l (in nanometers). The formal unit for specific
rotation is deg cm2 g–1 but scientific literature uses just degrees.
3.7 laurent’s half-shade PolariMeter LO6

It consists of two nicol prisms P1 and P2 as shown in Fig. 3.17. The nicol prism P1 acts as a polariser and
P2 acts as an analyser. The half-wave plate H is placed between nicol prism P1 and P2. This plate consists of
Polarisation 139
Monochromatic P1 P2
T
Source of
Light Telescope
S
H
Polariser Nicol Prism Polarimeter Tube Analyser Nicol Prism
(Filled with solution)
Figure 3.17
Y Y
Q¢ Q E-ray
P Q
q q
q q
X¢ X O-ray
C X¢ N¢ C N X
Q Q¢
Y´
Y¢
(a) (b)
Figure 3.18
two semi-circular plates, one half of which is made of glass and the other half of quartz and both halves are
attached together as shown in Fig. 3.18a. The thickness of the quartz plate is kept such that it introduces a
phase difference of p between the ordinary and extraordinary vibrations.
The monochromatic light from the source S is incident on convex lens, from which it emerges as parallel
beam and falls on the polarising nicol prism P1. The light emerging from the polariser P1 is plane-polarised
and falls on the half-shade plate H and the on a polarimeter tube T filled with optically active solution. Finally
the light emerged from tube falls on the analyser. Then we see emergent light with the help of a telescope.
The analysing nicol prism can be rotated about its axis. Its rotations are measured in term of angle q by using
circular scale.
3.8 biqUartz pOLariMeter LO6

The arrangement of biquartz polarimeter is the same as that of a Laurent’s half-shade polarimeter. However,
there are some different features of this polarimeter. For examples, it consists of a biquartz plate in place of
Laurent’s half-shade plate and it uses white light in place of monochromatic light.
A biquartz plate consists of two semicircular plates one of left-handed and the other of right-handed quartz,
both cut perpendicular to optic axis and joined together so as to form a complete circular plate. The left-
handed and right-handed quartz plates rotate the plane of polarisation of the incident beam in anti-clockwise
and clockwise directions, respectively. The thickness of both half-plates is the same and is adjusted such that
each rotates the plane of polarisation of yellow light through 90°. When a beam of white light is incident on
the polarising nicol P1 and then polarised light is incident on the X

Biquartz plate, it rotates the different wavelengths of white light to L R
different extents. For example, the red light colour is rotated the
r r
least, the yellow light is rotated through 90° and the violet light
rotates through a maximum angle.
In Fig. 3.19 XCX ¢ is the plane of polarisation of the incident plane-
polarised white light. Here, the analyser is set in such a way that its
Y Y´
principal section is parallel to XCX ¢. In this figure, YCY ¢ represents C
the plane of polarisation of yellow colour. Since with the use of
biquartz plate, the plane of polarisation of yellow colour YCY ¢ is
rotated through 90°, the vibration of yellow colour will be parallel
to XCX ¢. In this position, the field of view (both the halves) will b b
appear grey which is the combination of blue and red colours.
This position is called as tint of passage. If the analyser is rotated X´
slightly from the setting of tint of passage, one-half of the field of Figure 3.19
view appears blue and the other half-red.
3.9 SacchariMeter LO6

It is a type of polarimeter, which is used to determine the concentration of a sugar solution by measuring
the angle of rotation of the plane of polarisation of the polarised light passing through a tube containing the
solution. The concentration of the solution (m/V) is calculated by the formula obtained from the concept of
specific rotation
m q
=
V S ¥l
3.9.1 laurent saccharimeter
It consists of two Nicol prisms namely N1 and N2, as shown in Fig. 3.20a. Nicol N1 is used to polarise the
light, so it works as a polariser. Nicol N2 works as an analyser. If N1 and N2 are kept parallel, the light can pass
through them. Then these are said to be parallel. If N1 and N2 are not kept parallel so the light does not pass
through them, then the Nicols are said to be crossed. Some substances like quartz, sugar solution, etc. possess
the property of rotating the plane of polarisation of the light. The amount of this rotation can be measured by
determining the angle through which Nicol N2 is turned.
In order to determine the angle by which the Nicol N2 is rotated, we keep a circular sheet made of quartz and
glass just in front of polarising Nicol (Fig. 3.20b). Glass is of such thickness that it absorbs the same amount
of light as the quartz does. Light gets separated into two components when it just reaches the quartz plate.
These components pass through the plate with different velocities. Let these components be represented by
OP and OE when vibrations at O take place in the direction OA (Fig. 3.20b). There will be a gradual change
of phase between these components due to different transmission velocities. After a time this disturbance
will reach a point in the plate where component displacement is along OP and along OE ¢. The resultant
displacement is OA ¢. This difference is one half a period on leaving the quartz plate and the plate is said to
be half-wave plate. Light passes through the glass undisturbed and its oscillations are still along the direction
ED which is parallel to OA.
Polarisation 141
Light Lens Nicol N1 Tube Nicol N2 Telescope

Source
(a)
F A C
A¢ P A
E¢ E
O
D B
(b)
Figure 3.20
If Nicol N2 is kept with its short diagonal at right angle to OA ¢, this OA ¢ component is not transmitted while
OA passes through and the glass side appears illuminated. In fact, the light passes through from both sides
of the plate, but both sides are not equally illuminated. When both sides present the same illumination, the
principle plane of the Nicol is either along AB or normal to it. If both halves are equally dark, the Nicol is so
placed that the smaller components are transmitted. If the Nicol is set for equal illumination on both sides
and an active substance is interposed it will be necessary to rotate the Nicol to find the position of equal
intensities. In this case, the amount of rotation determines the angle of rotation of the plane of polarisation.
3.10 phOtOeLaSticity LO7

Photoelasticity is a phenomenon in which transparent isotropic substances such as plastic or glass could be
made optically anisotropic by applying mechanical stress. It is an experimental technique for measuring and
visualising stresses in structures by means of a phenomenon known as birefringence. The material becomes
birefringent under the compression and tension. This induced birefringence in the material is proportional
to the stress. It is found that the optic axis is induced in the direction of the stress in both the cases of
compression and tension. If nonuniform stress is applied to the sample, the phenomena of birefringence or
the retardation will not take place on a transmitted wave.
As discussed, two crossed polarisers ordinarily do not transmit light but when a photoelastic material is
placed between them and the plane of polarisation of light is not parallel to the principal axis of the stress,
some light will be transmitted. Under this situation, a complicated coloured fringe pattern is obtained, which
is helpful to measure the internal mechanical stresses in the samples.
When a beam of plane-polarised light is passed through a specimen of birefringent material or photoelastic
material, then the light is resolved along the two principal stress directions. Each component of the light
experiences different refractive indices. Then relative phase retardation (retardance) at any point on the
specimen material is directly proportional to the principal phase difference, i.e.,
R µ (s1 – s2)
R = Cd(s1 – s2)
Here R is the induced retardation, C is the stress optic coefficient and d is the thickness of the specimen
material, s1 and s2 are the orthogonal principal stresses. For example, if we consider the specimen to be a
plate then s1 would be the maximum principal stress in the vertical direction and s2 would be the minimum
principal stress in the horizontal direction.
The phenomenon of interference takes place when these two beams are brought together in a polariscope.
Then a colourful fringe pattern is obtained which depends on the retardance. Thus, study of these colourful
fringes provides the state of stress at the various points in the material. The locii of all the points on the
specimen for which (s1 – s2) remains constant under white light illumination are knows isochromatic region
and each such region corresponds to the definite colour. If the plane of polarisation of light is parallel to the
principal stress axis, then the wave will pass undeflected though the sample regardless of wavelength. Either
under normal condition or under stress, many ordinary materials show birefringence. For example, when
a crumpled piece of cellophane is introduced between cross polarisers, then it shows a striking variety of
colours.
s UmmarY
Following important points can be noted related to the matter presented in this chapter.
✦ Starting with a general introduction of light as an electromagnetic wave, the concept of polarisation of
light was introduced.
✦ Based on a mechanical experiment using a string, various types of polarisation of wave viz. linearly
polarised wave, circularly polarised wave and elliptically polarised wave were discussed.
✦ Along with the appropriate figures, the difference between unpolarised light and polarised light was
made clear.
✦ The important features of polarised light including the direction of polarisation, plane of polarisation
and plane of vibration were discussed.
✦ Various means of production of plane-polarised light were introduced.
✦ Polaroid filter is commonly used for generating the plane-polarised light. So it was discussed in detail
including the alignment of its long chain molecules and the polarisation axis.
✦ Malus’ law of intensity was discussed related to the intensity of light emerging from a Polaroid filter
(analyser) whose transmission axis make an angle with the axis of another Polaroid which is used to
produce plane-polarised light.
✦ Concept of Brewster’s angle was introduced for which the reflected light can be obtained as fully
polarised light.
✦ Biot’s polariscope was discussed in short. This can be used for producing and detecting the plane-
polarised light by the reflection.
✦ Polarisation can also take place by the reflection of light which occurs when a beam of light passes
from one material into another material. When the unpolarised light is passed through a particular type
of crystal, the light gets split into two rays. These waves are polarised parallel and perpendicular to a
particular direction in the crystal. This particular direction is called as optic axis. So optic axis was
discussed in detail.
Polarisation 143
✦ The principal section of a crystal was also discussed.

✦ Geometry of calcite crystal was discussed with relation to the double refraction which is nothing but the
splitting of light wave into two rays in the crystal. These two rays are referred to as ordinary ray (O-ray)
and extraordinary ray (E-ray).
✦ Experimental observation of polarisation by double reflection was discussed in detail and an explanation
was also given.
✦ Huygens’ theory of double refraction was discussed qualitatively.
✦ Out of the two rays (E-ray and O-ray) obtained by double refraction, one ray can be suppressed with the
use of Nicol prism. So it was discussed in detail.
✦ Since E-ray and O-ray are plane-polarised waves, their superposition can lead to the resultant wave as
circularly polarised wave or elliptically polarised wave. In view of this, the concept of quarter-wave
plate and half-wave plate for the inclusion of phase change of p/2 or p between E-ray and O-ray was
given.
✦ Theory of production of plane-polarised light, circularly polarised light and elliptically polarised light
was discussed in detail.
✦ There are some substances that rotate the plane of polarisation of light which is passed through them.
These substances are called optically active substances and this property is called as optical activity. So
this was discussed along with the inclusion of specific rotation.
✦ Polarimeters were discussed which are used to produce the polarised light and also to determine the
rotation of its plane of polarisation when passed through an optically active substance. Specifically
Laurent’s half shade polarimeter and biquartz polarimeter were discussed.
✦ The concept of polarimeter was extended to Saccharimeter, which is used to determine the concentration
of sugar solution based on the measurement of angle of rotation of the plane of polarisation of light.
✦ Lastly photoelasticity was discussed which is a phenomenon in which transparent isotropic substances
such as plastic or glass could be made optically anisotropic by applying mechanical stress. It is an
experimental technique for measuring and visualising stress in structures by means of a phenomenon
known as birefringence. The material becomes birefringence under the compression and tension.
s olved e XamPles
E xamplE 1 Refractive index of glass is 1.5. Calculate Brewster’s angle for it. Also calculate the angle of
refraction.
Solution Given m = 1.5
Brewster’s law says, m = tan ip
3
or tan ip = m = 1.5 =
2
ip = 56.31°
r = 90 – ip = 90 – 56.31 = 33.69°
r = 33.69°
E xamplE 2 The refractive index for water is 1.33. Calculate the polarising angle for water.
Solution Given m = 1.33
Formula used is tan ip = m = 1.33
ip = tan–1 (1.33) = 53.06°
ip = 53.06°
E xamplE 3 The refractive indices of glass and water are 1.54 and 1.33, respectively. For which case polarising
angle will be greater: for a beam incident from water to glass or for a beam incident from glass to water?
Solution Given mg = 1.54 and mw = 1.33
Formula used is tan ip = m
mg 1.54
For water to glass wmg = = = 1.16
mw 1.33
So ip = tan–1 (wmg) = tan–1(1.16)
ip = 49.23°
m w 1.33
For glass to water gmw = = = 0.864
m g 1.54
So ip = tan–1(0.864) = 40.82°
Hence, polarising angle (ip) is greater for a beam incident from water to glass.
E xamplE 4 If the polarising angle of a piece of glass for green light is 60°, calculate the angle of minimum
deviation for a 60° prism made of same glass.
Solution Given ip = 60°
m = tan ip = tan 60°
or m = 1.732
È A + dm˘
sin Í
Î 2 ˚˙ ,
In case of prism m = where dm is the angle of minimum deviation and A is the prism angle
A
sin
2
Here, A = 60° and m = 1.732 dm = ?
È 60∞ + d m ˘
sin Í
\ Î 2 ˚˙ = 1.732
60∞
sin
2
È 60∞ + d m ˘ 1
or sin Í
Î 2 ˙˚ = 1.732 ¥ 2 = 0.866
60∞ + d m
= 60∞
2
or dm = 60°
Polarisation 145
E xamplE 5 Determine the Brewster’s angle for a glass of refractive index 1.5 when it is immersed in water
of refractive index 1.33.
Solution Given mg = 1.5 and mw = 1.33
1.5
Therefore, the refractive index of glass w.r.t water = = 1.128
1.33
By Brewster’s law, tan ip = m, where ip = Brewster’s angle
ip = tan–1m = tan–1(1.128)
ip = 48.4°
E xamplE 6 A ray of light is incident on a glass plate of refractive index 1.732 at a polarising angle. Find the
angle of incidence and angle of refraction.
Solution Given m = 1.732, tan ip = m, where ip = angle of polarisation = angle of incidence
ip = tan–1(m) = tan–1(1.732) = 60°
Now ip + r = 90° or r = 90° – 60° = 30
So the angle of incidence is 60° and angle of refraction is 30°.
E xamplE 7 If the angle between a polariser and analyser is 60°, what will be the intensity of transmitted light
for original intensity of incident light as I0?
Solution Given q = 60°
According to Malus’ law I = I0 cos2q
I = I0 cos260 = I0(0.5)2
I = 0.25I0
E xamplE 8 Unpolarised light is incident on two polarising sheets placed one on top of the other. What must
be the angle between the characteristics direction of sheets if the intensity of transmitted light is (i) 1/3 of
maximum intensity of the transmitted beam and (ii) 1/3 of intensity of incident beam. Assume that sheet
reduces the intensity of unpolarised light by exactly 50%.
Solution Consider intensity of unpolarised light as I0. The intensity of polarised light transmitted by the first sheet would be
1
I= I0
2
Case-I: q1 = ?
1
I1 = I
3
1
I1 = I cos2q1 or I = I cos2q1
3
1
or cos q1 =
3
or q1 = 54.74°
1 2
Case-II: I2 = I0 = I
3 3
2
I2 = I cos2q2 or I = I cos2q2
3
2
or cos q2 = or q2 = 35.26°
3
E xamplE 9 Two Nicols have parallel polarising directions so that the intensity of transmitted light is maximum.
Through what angle must either Nicol be turned if intensity is to drop by one-fourth of its maximum value?
3 3
Solution The transmitted intensity will be of incident intensity, i.e., I0
4 4
Malus’ Law I = I0 cos2q
3
or I0 = I0 cos2q
4
3
cos q =
2
or q = 30°
E xamplE 10 Two Nicol prisms are so arranged that the amount of light transmitted through them is
maximum. What will be the percentage reduction in the intensity of the incident light when the analyser is
rotated through (i) 30°, (ii) 45° (iii) 60° and (iv) 90°?
Solution
I = I0 cos2q Malus Law (i)
I
cos 2 q = (ii)
I0
where I0 is the intensity of incident light and I is the intensity of transmitted light.
Therefore, percentage reduction in the intensity of incident light is
I0 - I Ê Iˆ
= ¥ 100 = Á1 - ˜ ¥ 100 (iii)
I0 Ë I 0¯
By using Eqs. (ii) and (iii), we get percentage reduction in intensity
= (1 – cos2q) ¥ 100
(i) For q = 30°, Percentage reduction in intensity
= (1 – cos230°) ¥ 100
= [1 – (0.866)2] ¥ 100
= 25%
(ii) For q = 45°, %reduction in intensity
= (1 – cos245°) ¥ 100 = 50%
(iii) For q = 60°, % reduction in intensity
= (1 – cos260) = 75%
and,
(iv) For q = 90°, % reduction in intensity
= (1 – cos290) = 100%
E xamplE 11 Two polaroids are adjusted so as to obtain maximum intensity. Through what angle should one
polaroid be rotated to reduce the intensity to (i) half (ii) one fourth.
Polarisation 147
Solution Formula used is I = I0cos2q (i)

I
or = cos2q (ii)
I0
I1 1
Case (i) =
I0 2
1
Then, by Eq. (i), we have cos2q =
2
1
or cos q =
2
or q = 45°
I1 1
Case (ii) =
I0 4
1
\ cos2q = [using eqn (i)]
4
1
cos q =
2
q = 60°
E xamplE 12 Calculate the thickness of a half-wave plate of a quartz for a wavelength of 5000 Å. Here
me = 1.553 and m0 = 1.544
Solution Given l = 5000 Å, me = 1.553 and m0 = 1.544
l 5.0 ¥ 10-7
For Half-wave plate t= = = 2.78 ¥ 10–5 m.
2( me - m0 ) 2 ¥ (1.553 - 1.544)
E xamplE 13 Calculate the thickness of quarter-wave plate for light of wavelength 5893 Å, given refractive
indices for ordinary ray and extraordinary ray as 1.544 and 1.533, respectively.
Solution Given mo = 1.554, me = 1.533 and l = 5.893 ¥ 10–7 m
l 5.893 ¥ 10-7
Formula used is t= = = 70.15 ¥ 10–7
4( m0 - me ) 4 ¥ (1.554 - 1.533)
= 7.02 ¥ 10–6 m. = 7.02 mm
E xamplE 14 Calculate the thickness of a quarter-wave plate of wavelength 5890 Å.

(i) Given mo = 1.55 and me = 1.50 (ii) Given mo = 1.55 and me = 1.57
Solution Given l = 5.89 ¥ 10–7 m
l 5.89 ¥ 10-7
(i) Formula used is t = =
4( m0 - me ) 4 ¥ (1.55 - 1.50)
= 2.945 ¥ 10–6 m
t = 2.945 mm
l 5.89 ¥ 10-7
(ii) t = =
4( me - m0 ) 4 ¥ (1.57 - 1.55)
= 7.36 ¥ 10–6 m
or t = 7.36 mm
E xamplE 15 Calculate the thickness of a calcite plate which would convert plane-polarised light into
circularly-polarised light. The principal refractive indices are mo = 1.658 and me = 1.486 at the wavelength
5890Å of light used.
Solution Given mo = 1.658, me = 1.486 and l = 5.89 ¥ 10–7 m
To convert plane-polarised light into circularly polarised light, path difference must be l/4.
l 5.89 ¥ 10-7
Hence, t= =
4( m0 - me ) 4 ¥ (1.658 - 1.486)
= 8.56 ¥ 10–7 m
E xamplE 16 Plane-polarised light passes through a quartz plate with its optic axis parallel to the faces.
Calculate the least thickness of the plate for which the emergent beam will be plane-polarised. (Given
me = 1.5533, mo = 1.5442 and l = 5 ¥ 10–5 cm).
Solution Given me = 1.5533, m0 = 1.5442, and l = 5 ¥ 10–5 cm
l 5.0 ¥ 10-7
t= =
2( me - m0 ) 4 ¥ (1.5533 - 1.5442)
= 1.37 10 5 m
E xamplE 17 Find the thickness of a quarter-wave plate when the wavelength of light is equal to 5890Å
m0 = 1.55 and me = 1.54.
Solution Given l = 5.89 ¥ 10–7 m, mo = 1.55 and me = 1.54
l 5.89 ¥ 10-7
Formula used is t= =
4( m0 - me ) 4(1.55 - 1.54)
t = 1.47 ¥ 10–5 m
E xamplE 18 Quartz has refractive indices 1.553 and 1.544. Calculate the thickness of the quarter-wave plate
for sodium light of wavelength 5890 Å.
Solution Given me = 1.553 and mo = 1.544
l = 5.89 ¥ 10–7 m
5.89 ¥ 10-7
t=
4 ¥ (1.553 - 1.544)
t = 1.63 ¥ 10–5m
E xamplE 19 Plane-polarised light (l = 5 ¥ 10–7 m) is incident on a quartz plate cut parallel to the optic axis.
Find the least thickness of the plate for which the ordinary and extraordinary rays combine to form a plane-
polarised light on emergence. What multiples of this thickness would give the same result? The indices of the
refraction of quartz are me = 1.5533 and m0 = 1.5542.
Solution Given mo = 1.5442, me = 1.5533 and l = 5 ¥ 10–7 m
In the given case the quartz plate must act as a half-wave (l/2) plate, then formula used is
l 5 ¥ 10-7
t= =
2( me - m0 ) 2(1.5533 - 1.5442)
Polarisation 149
or t = 2.75 ¥ 10–5 m
The thickness that would give the similar result should be t, 3t, 5t ... i.e., 2.75 ¥ 10–5 m, 8.25 ¥ 10–5 and so on.
E xamplE 20 On introducing a polarimeter tube 25 cm long and containing sugar solution of unknown
strength, it is found that the plane of polarisation is rotated through 10°. Find the strength of the sugar solution
in g/cm3 (Given that the specific rotation of sugar solution is 60° per decimeter per unit concentration).
Solution Given q = 10°, S = 60° and l = 25 cm
q
Formula used is s = where l is in decimeter
lc
10q
or s= for l in cm
lc
10q 10 ¥ 10
c= =
or ls 25 ¥ 60
or c = 0.067 g/cc
E xamplE 21 Compute the specific rotation if the plane of polarisation is turned through 26.4°, traversing
20 cm length of 20% sugar solution.
Solution Given q = 26.4°, l = 20 cm and c = 20% = 0.2 g/cm3
10q 10 ¥ 26.4
Formula used is S= =
lc 20 ¥ 0.2
S = 66°
E xamplE 22 The plane of polarisation of plane-polarised light is rotated through 6.5° in passing through
a length of 2.0 decimeter of sugar solution of 5% concentration. Calculate the specific rotation of the sugar
solution.
Solution Given q = 6.5°, l = 2 dm and c = 5% = 0.05 g/cc
q 6.5
Formula used is S= =
lc 2 ¥ 0.05
= 65° (dm)–1 (gm/cc)–1
E xamplE 23 80 gm of impure sugar when dissolved in a litre of water gives an optical rotation of 9.9° when
placed in a tube of length 20 cm. If the specific rotation of sugar is 66°, find the percentage purity of the sugar
sample.
Solution Given q = 9.9°, l = 20 cm = 2.0 dm and S = 66°
q q
Formula used S= or c =
lc ls
9.9
or c= = 0.075 g/cc
2.0 ¥ 66
= 75 gm/L [1 litre = 103 cc]
80 g impure sugar is dissolved in one litre of water in which 75 g sugar is pure. Therefore, percentage of pure sugar is
75
¥ 100 = 93.75%
80
E xamplE 24 A 20 cm long tube containing sugar solution rotates the plane of polarisation by 11°. If the
specific rotation of sugar is 66°, determine the strength of the solution.
Solution Given q = 11°, l = 20 cm and S = 66°
10q 10q 10 ¥ 11
Formula used is S= or c = =
lc lS 20 ¥ 66
= 0.0833 g/cm3
E xamplE 25 Calculate the specific rotation if the plane of polarisation is turned through 26.4°, traversing
20 cm length of 20% sugar solution.
Solution Given q = 26.4°, l = 20 cm, ad c = 20% = 0.2 g/cc
10q 10 ¥ 26.4
Since l is given in cm, specific rotation S = = = 66°
lc 20 ¥ 0.2
E xamplE 26 A sugar solution in a tube of length 20 cm produces optical rotation of 13°. The solution is
then diluted to one-third of its previous concentration. Find the optical rotation produced by 30 cm long tube
containing the diluted solution.
c
Solution l = 20 cm, q = 13°, c ¢ = and l ¢ = 30 cm
3
10 ¥ q
Formula used is S=
lc
10q 10q ¢
Therefore, S= =
lc l ¢c¢
Ê l ¢ ˆ Ê c¢ ˆ Ê 30 ˆ Ê c /3 ˆ
q¢ = Á ˜ Á ˜ q = Á ˜ Á 13∞
Ë l ¯Ë c¯ Ë 20 ¯ Ë c ˜¯
3 1
= ¥ ¥ 13∞ = q ¢ = 6.5°
2 3
o bJective t YPe Q Uestions
Q.1 Which of the following phenomenon tells about the transverse nature of light waves?
(a) interference (b) diffraction (c) polarisation (d) photoelectric effect
Q.2 Plane-polarised light has vibrations
(a) in one direction perpendicular to the direction of propagation
(b) along the direction of propagation
(c) in all directions perpendicular to the direction of propagation
(d) in two directions perpendicular to the direction of propagation
Q.3 Polarised light can be produced by
(a) reflection (b) refraction (c) double refraction (d) all of them
Q.4 At polarising angle, the reflected and refracted rays make angle
(a) 90° (b) 180° (c) 30° (d) none of these
Q.5 Brewster’s law in terms of refractive index can be expressed as
(a) m = sin ip (b) m = cos ip (c) m = tan ip (d) m = cot ip
Polarisation 151
Q.6 According to the Malus’ Law, the intensity of polarised light emerging through the analyser varies as
I 2
(a) I0 cos2q (b) I0 sin2q (c) I0 cos q (d) 0 cos q
2
Q.7 Malus’ discovered a simplest method for polarisation of light by reflection in the year
(a) 1808 (b) 1908 (c) 1856 (d) none of these
Q.8 Which one is the example of uniaxial crystal?
(a) calcite (b) tourmaline (c) quartz (d) all of them
Q.9 Which one is the example of biaxial crystal?
(a) sodium chloride (b) tourmaline (c) aragonite (d) none of them
Q.10 What happens if the ordinary unpolarised light is passed through a uniaxial crystal?
(a) light is split into two rays (b) light remains unaffected
(c) light is split into more than two rays (d) none of them
Q.11 What happens to O and E-rays if they travel along the optic axis?
(a) both rays travel with same velocity (b) O-ray travels faster than E-ray
(c) E-ray travels faster than O-ray (d) none of them.
Q.12 How many principal sections of uniaxial crystal has?
(a) 6 (b) 3 (c) 5 (d) 2
Q.13 At what angle of incidence of plane-polarised light with quarter-wave plate elliptically polarised light
becomes circularly polarised?
(a) 90° (b) 45° (c) 60° (d) 30°
Q.14 How much phase change is introduced by a quarter-wave plate between ordinary and extraordinary
rays?
(a) p (b) 2p (c) p/2 (d) p/4
Q.15 Dextrorotatory optically active substance rotates the plane of vibrations
(a) in clockwise direction (b) in anti-clockwise direction
(c) by 180° (d) none of them
Q.16 Which of the following relation is true for quartz crystal?
(a) me > mo (b) mo > me (c) mo = me (d) none of these
Q.17 Which of the following relation is true for quartz crystal?
(a) vo > ve (b) ve > vo (c) ve = vo (d) none of these
Q.18 The substance that is capable of rotation of plane of vibration is known as
(a) optically active (b) optically inactive (c) both (a) and (b) (d) none of these
Q.19 If two polarising and analysing nicols are at 90° than, the emergent light passed through the analysing
nicol becomes
(a) maximum (b) minimum (c) zero (d) none of these
s Hort -a nsWer Q Uestions
Q.1 What do you understand by polarisation?

Q.2 Experimentally show that only the transverse waves can be polarised.
Q.3 What is the difference between polarised and unpolarised light?
Q.4 Distinguish between plane of vibration and plane of polarisation.

Q.5 What is Brewster’s law?
Q.6 State Malus’ law?
Q.7 What is the relation between angle of polarisation and refractive index?
Q.8 What do you understand by double refraction?
Q.9 What are ordinary and extraordinary rays?
Q.10 What is optic axis?
Q.11 What is principal section of the crystal?
Q.12 Differentiate between uniaxial and biaxial crystals?
Q.13 What is Nicol prism?
Q.14 Distinguish between linearly, circularly and elliptically polarised light.
Q.15 What are quarter-wave and half-wave plates?
Q.16 What do you mean by optical activity?
Q.17 What do you understand by specific rotation?
Q.18 What do you understand by photoelasticity?
Q.19 Write a short note on
(a) Polarised and unpolarised light,
(b) Double refraction,
(c) Nicol prism,
(d) Photoelasticity,
(e) Quarter and half-wave plate,
(f) Biquartz polarimeter,
(g) Half-shade polarimeter.
P ractice P roblems
Q.1 Explain polarisation on the basis of electromagnetic theory of light.

Q.2 Explain the method and cause of production of plane-polarised light by reflection.
Q.3 State and derive Brewster’s law. What does the law become when the rays of light travel from denser
medium to rarer medium?
Q.4 By Brewster’s law show that light incident on a transparent substance at polarising angle gives reflected
and refracted rays at right angle to each other.
Q.5 How can you explain the transverse nature of a wave with the help of phenomenon of polarisation
using mechanical and light experiments?
Q.6 Distinguish between polarised light and unpolarised light. Also define plane of vibration and plane of
polarisation.
Q.7 How can you obtain linearly polarised light from a beam of a non-polarised light?
Q.8 How would you obtain plane-polarised light? Name various techniques and discuss one out of them
which is based on double refraction phenomenon.
Q.9 What do you understand by double refraction? Explain how would you use the phenomenon to produce
plane-polarised light and circularly polarised light?
Polarisation 153
Q.10 How do you use the phenomenon of double refraction to produce a plane-polarised light? Explain in
detail.
Q.11 (a) Explain the phenomenon of double refraction in uniaxial crystal. (b) What are quarter-wave and
half-wave plates? Explain their use in the study of different types of polarised light.
Q.12 Explain the phenomenon of polarisation of light. Describe the construction of a Nicol Prism, and show
how it can be used as a polariser and as an analyser.
Q.13 Discuss the principle, construction and working of Nicol prism as polarised.
Q.14 Draw diagrams and discuss double refraction through uniaxial crystals due to a plane wave when
(i) optic axis is inclined to the upper face but lying in the plane of incidence.
(ii) optic axis is parallel to the upper face but lying in the plane of incidence.
(iii) optic axis is parallel to the upper face but perpendicular to the plane of incidence.
(iv) optic axis is perpendicular to the upper face.
Q.15 What do you understand by optical rotation? Explain Fresnel’s theory of the rotation of the plane of
polarisation. How would you increase the sensitivity of a pair of crossed Nicols?
Q.16 How would you distinguish between circularly polarised light and unpolarised light?
Q.17 How would you distinguish between plane, circularly and elliptically polarised light?
Q.18 What are plane-polarised circularly polarised and elliptically polarised light? Explain their production
with the help of mathematical equations. Give the salient features of biquartz device.
Q.19 Give two differences between Laurent’s half shade polarimeter and biquartz polarimeter.
Q.20 Define specific rotation. Describe the construction and working of Laurent’s half-shade polarimeter.
Discuss the relative merits of biquartz polarimeter and half-shade polarimeter.
Q.21 What is optical activity? Describe the construction, theory and working of biquartz polarimeter to find
the optical rotation of a solution and also discuss the action of biquartz plate in it.
Q.22 What do you understand by a half and quarter-wave plate? Give the theory and construction of Laurent’s
half-shade polarimeter.
Q.23 What is specific rotation? Describe the construction and working of biquartz polarimeter to find the
specific rotation of sugar solution and discuss the utility of biquartz plate in it.
Q.1 A beam of light is incident on a glass plate at an angle of 58°6 ¢ and the reflected beam is completely
plane-polarised. Find the refractive index of glass. [Ans: 1.6]
Q.2 Refractive index of water is 1.33. Calculate the angle of polarisation for light reflected from the surface
of a pond. [Ans: 53.06°]
Q.3 Critical angle for refraction for glass to air is 40°. Calculate the polarising angle for glass.
[Ans: 57.3°]
Q.4 A beam of light traveling in water strikes a glass plate which is also immersed in water. When the angle
of incidence is 51°, the reflected beam is found to be plane-polarised. Calculate the refractive index of
glass. [Ans: 1.235]
Q.5 A polariser and an analyser are set in such a way that the intensity of the emergent light is maximum.
What percentage of the maximum intensity of light is transmitted from the analyser if either is rotated
by 30°, 45° and 60°? [Ans: 75%, 50%, 25%]
Q.6 Two Nicols have parallel polarising directions so that the intensity of transmitted light is maximum.
Through what angle must either Nicol be turned if the intensity is to drop by one fourth of its maximum
value? [Ans: 30°]
Q.7 An analysing Nicol examines two adjacent plane-polarised beams A and B whose planes of polarisaton
are mutually perpendicular. In one postion of the analyser, beam B shows zero intensity. From this
position a rotation of 30° shows the two beams as matched (i.e., of equal intensity). Deduce the
intensity ratio IA/IB of the two beams [Ans: 1/3]
Q.8 Find the thickness of a quarter-wave plate for the wavelength 5890Å of light, when m0 = 1.55 and
me = 1.54. [Ans: 1.4725 ¥ 10–5 m]
Q.9 Find the thickness of calcite plate which would convert plane-polarised light into circularly polarised
light. The refractive indices are m0 = 1.658 and me = 1.486 at the wavelength of light used as 5890Å.
[Ans: 8.56 ¥ 10–9 m]
Q.10 Calculate the thickness of a quarter-wave plate of quartz for sodium light of wavelength 5893Å. The
refractive indices of quartz for ordinary and extra-ordinary rays are 1.5442 and 1.5533 respectively.
[Ans: 1.61 ¥ 10–5 m]
Q.11 Calculate the thickness of a doubly refracting crystal plate required to introduce a path difference of
l/2 between the ordinary and extraordinary rays when l = 6000Å, m0 = 1.55 and me = 1.54.
[Ans: 5 ¥ 10–5 m]
Q.12 Calculate the thickness of (i) a quarter-wave plate and (ii) a half-wave plate. Given that me = 1.553 and
m0 = 1.544 and l = 5000Å. [Ans: (i) 1.39 ¥ 10–5 m (ii) 2.78 ¥ 10–5 m]
Q.13 A plane-polarised light is incident on a piece of quartz cut parallel to the axis. Find the least thickness
for which the ordinary and extraordinary rays combine to form plane-polarised light. Given m0 =
1.5442, me = 1.55 and l = 5 ¥ 10–5 cm. [Ans: 2.75 ¥ 10–4 m]
Q.14 For calcite m0 = 1.658 and me = 1.486 for sodium light. Calculate the minimum thickness of the quarter-
wave plate for calcite. [Ans: t = 8.56 ¥ 10–7 m]
Q.15 A 20 cm long tube containing sugar solution rotates the plane of polarisation by 11°. If the specific
rotation of sugar is 66°, calculate the strength of the solution. [Ans: 0.0833 g/cc]
Q.16 A 200 mm long tube containing 48 cm3 of sugar solution produces an optical rotation of 11° when
placed in a saccharimeter. If the specific rotation of sugar solution is 66°, calculate the quantity of
sugar contained in the tube in the form of a solution. [Ans: 4.0 g]
Q.17 A 20 cm long tube is filled with a solution of 15 g of cane sugar in 100 cc of water. Find the angle
of rotation of the plane of polarisation of a beam of plane-polarised light when it passes through the
solution. Specific rotation for cane sugar = 66.50 per dm/g per cm3. [Ans: 20°]
Lasers and Holography 4
Learning Objectives
LO1 Learn about absorption of radiation LO4 Discuss the application of laser and
and different types of emissions laser cooling
LO2 Understand the phenomenon LO5 Explain holography versus conventional
of population inversion and photography, recording and
characteristics of laser light reconstruction of image on a holograph
LO3 Know about the components and types LO6 Illustrate types of holograms
of lasers. LO7 Evaluate the applications of holography
Introduction
In the previous chapters, interesting phenomena of interference and diffraction of light including its
polarisation have been investigated in detail. It was discussed that the interference has scientific as well
as engineering applications. The concept of interference is applied to testing the surface quality of optical
components and this led to the development of flatness interferometers. An exciting use of the concept
of interference is made in the preparation of nonreflecting or antireflecting coatings that are applied to
surfaces of lenses (for example, eye glass lenses) and other optical devices for reducing the reflections and
hence in improving the efficiency of the system like telescope. However, you would have learnt that in
order to realise the above mentioned phenomena in an efficient way there is a need of using the coherent
and monochromatic sources as the phase of incoherent source (light) varies randomly with time and
position. This need of monochromatic and coherent sources contributed to the birth of a special type of
device that amplifies light and produces a highly intense and highly directional beam which mostly has a
very pure wavelength. This device is called LASER. Lasers are available with power ranging roughly from
1 nW (= 10–9 W) to 105 PW (1 PW = 10–15 W) and with frequency ranging from 100 GHz (1 GHz =
109 Hz) to 100 PHz. Nowadays the lasers with pulse duration as short as ~ 1 fs (= 10–15 s) are available
with their pulse energies as high as 10 kJ.
The name LASER is an acronym of Light Amplification by Stimulated Emission of Radiation. The immediate
originator to the LASER is the MASER, formerly acronym of Microwave Amplification by Stimulated
Emission of Radiation. Since the techniques have been extended to the infrared and optical regions, it has
now come to stand for Molecular rather than Microwave amplification. A laser uses some processes that
amplify light signals. These processes mainly include stimulated emission and optical feedback provided by
mirrors. The stimulated emission takes place in amplifying medium contained by the laser. The application
of set of mirrors is to feed the light back to the amplifying medium so that the developed beam is
grown continuously. The key concept for realisation of the laser operation is the principle of coherence
accompanying stimulated emission.
4.1 AbsOrptiOn And EmissiOn Of rAdiAtiOn LO1

It is well known that an atom can be excited by supplying energy with an amount equal to the difference of its
any two energy levels. Then after a very short duration of time the atom shall radiate energy when it comes
down to its lower energy state. An electron undergoes a transition between two energy states E1 and E2 if
the atom emits or absorbs a photon of appropriate energy as per the relation E2 – E1 = hn, where h is Planck
constant and n is the frequency of radiation.
4.1.1 Absorption of radiation
At low temperatures, most of the atoms stay in lower energy states. 2 E2 2 E2
If an atom is initially in the lower energy state E1, it can be raised to hv
the higher energy state E2 by the absorption of a photon of energy hn, 1 E1 1 E1

Before After
as shown in Fig. 4.1. This is know as absorption of radiation and is
Figure 4.1
represented by
E2 = E1 + hn
fi E2 – E1 = DE = hn (i)
The probability of occurrence of this absorption from state 1 to state 2 is proportional to the energy density
u(n) of the radiation
P12 = B12 u(n) (ii)
Where the proportionality constant B12 is known as the Einstein’s coefficient of absorption of radiation.
Take an example of electron transition associated with visible and ultraviolet radiation interactions with
matter. Here the absorption of a photon occurs only when the quantum energy of the photon precisely matches
the energy gap between the initial and final states. In such interaction of radiation with matter, if there is no
pair of energy state such that the photon energy can elevate the system from the lower to upper state, then the
matter will be transparent to that radiation.
4.1.2 spontaneous Emission 2 E2 2 E2
hv
If an atom is initially in the upper state E2, it can come down to lower 1 E1 1 E1
state E1 by emitting a photon of energy hn as shown in Fig. 4.2. Before After
This is known as spontaneous emission. This is the natural radiation Figure 4.2
decay process that is inherent in all excited states of all materials. However, such emission is not always the
dominant decay process.
The probability of occurrence of this spontaneous emission transition from state 2 to state 1 depends only on
the properties of states 2 and 1 and is given by
P12¢ = A21
(iii)
Where A21 is known as the Einstein’s coefficient of spontaneous emission of radiation.
4.1.3 Stimulated (Induced) Emission

Einstein was the first to point out a third possibility of induced 2 E2 2 E2
hv
emission, in which an incident photon of energy hn causes a hv
hv
transition from upper state E2 to the lower state E1, as shown in 1 E1 1 E1

Before After
Fig. 4.3. This occurs when
Figure 4.3
hn = DE = E2 – E1
In a system of atoms in thermal equilibrium, the number of atoms in the ground state is generally much
greater than in a higher energy state. This is known as normal population of atoms among the available
energy states. A state in which the number of atoms in higher energy state is greater than that of lower energy
state is known population inversion.
Therefore, the incoming photon stimulates the transition to the lower state and produces a second photon of
the same energy, when a sizable population of electrons resides in upper level (Fig. 4.4). In this condition it
is called a population inversion. This population inversion sets the stage for stimulated emission of multiple
photons. This is the precondition for the light amplification in a laser. Since the emitted photons have a
definite time and phase relation to each other, the light has a high degree of coherence. If these emitted
photons are passed through an assembly of atoms, which fulfil the condition of population inversion, these
are amplified. This amplification is very much clear from Fig. 4.4, which shows multiplication of photons
emitted during the process of stimulated emission.
(energy of each photon is hv)
E2 E2
hv
hv
hv hv hv
hv
hv
hv
(energy of each
hv hv
hv
hv hv photon is hv)
hv
E1 E1 hv
Figure 4.4
The probability of occurrence of stimulated emission transition from the upper level 2 to the lower level 1 is
proportional to the energy density u(n) of the radiation and is expressed as
P21¢¢ = B21 u (v) (iv)
Where B21 is the Einstein’s coefficient of stimulated emission of radiation.
Thus, the total probability of emission transition from the upper level 2 to the lower level 1 is given by
P21 = P21¢ + P21¢¢
or P21 = A21 + B21 u(n) (v)

4.1.4 Relation between Einstein’s Coefficients

Let N1 and N2 be the number of atoms at any instant in the state 1 and 2, respectively. The probability of
absorption transition for number of atoms from state 1 to 2 per unit time is given by
N1P12 = N1B12 u (n) (vi)
The total probability of transition for number of atoms from state 2 to 1, either by spontaneously or by
stimulated emission per unit time is given by
N2P21 = N2[A21 + B21 u(n)] (vii)
In thermal equilibrium at temperature T, the absorption and emission probabilities are equal and thus, we can
write
N1P12 = N2P21
or N1B12 u(n) = N2[A21 + B21 u(n)]
N 2 A21
or u (v ) =
N1B12 - N 2 B21
A21 1
or u (v ) = (viii)
B21 ( N1 / N 2 ) ( B12 / B21 ) - 1
But according to Einstein

B12 = B21 (ix)
Then from Eq. (viii) and (ix), we get
A21 1
u (v ) = (x)
B21 ( N1 / N 2 ) - 1
According to Boltzmann’s law, the distribution of atoms among the energy states E1 and E2 at the thermal
equilibrium at temperature T is given by
N1 e - E1 / kT
= = e( E2 - E1 )/ kT (xi)
N 2 e - E2 / kT
N1
or = e hv / kT (xii)
N2
Here k is the Boltzmann constant.
From Eq. (x), we can write
A21 1
u (v ) = hv / kT
(xiii)
B21 e -1
Planck’s radiation formula yields the energy density of radiation u(v) as
8p hv3 1
u (v ) = 3 hv / kT
(xiv)
c e -1
Comparing Eq. (xiii) and (xiv), we get

A21 8p hv3
= (xv)
B21 c3
Equation (xv) gives the relation between the probabilities of spontaneous and stimulated emissions. This is
also known as the relation between the Einstein’s coefficients A and B.
4.2 pOpuLAtiOn invErsiOn LO2

For a system with three energy states viz. E1 (population N1), E2 (population N2) and E3 (population N3)
in equilibrium such that E1 < E2 < E3, the uppermost level E3 is populated least whereas the lowest level
E1 is populated most. Since the population in these states follow the trend N1 > N2 > N3, the system shall
absorb photons rather than emitting them. However, when a sizable population of electrons is achieved in
the upper levels, the condition is known as population inversion (a non-equilibrium state). This condition
sets the stage for stimulated emission of radiation, i.e., multiple photons, as the first few randomly emitted
spontaneous photons trigger stimulated emission of more photons and those stimulated photons induce still
more stimulated emissions, and so on.
The population inversion is the precondition for the light amplification occurring in LASER. In order to
achieve this condition, a multilevel scheme is used. For example, the atoms are pumped into the highest level
of the three levels. Then spontaneous de-excitation from this pumped level to the metastable level takes place
and the laser emission occurs between the metastable level and the ground state. It is clear that energy has
to be supplied to the laser medium in order to raise atoms from the lower level to the excited level and for
maintaining population at the excited level at a value greater than that of the lower energy. One can think that
heating the material can solve this purpose but this only increases the average energy of the atoms and does
not enhance the population in the higher level. However, following schemes may be adopted to achieve the
population inversion.
4.2.1 Schemes for Population Inversion

To discuss schemes for the population inversion in various energy level systems, we will first prove that the
two-level system is not appropriate to achieve this condition of population inversion.
4.2.1.1 Two-level System
Consider the case of two-level system having energies E1 and E2 such that E2 > E1. We can easily find that
the Einstein coefficients (or constants) for the upward (B12) and downward (B21) transitions are equal, i.e.,
B12 = B21. It means, even with strong pumping, the population distribution in upper and lower levels can only
be made equal, i.e., the optical pumping will at most only achieve equal population of a two-level system.
This is due to the fact that the probabilities for raising an electron to the upper level and inducing the decay of
an electron to the lower level (stimulated emission) are exactly the same. In other words, we can say that the
numbers of electrons going up and coming down will be the same when both the levels are equally populated.
So, we cannot achieve population inversion in the case of two energy levels system. Therefore, optical as
well as any other pumping method needs either three or four level systems to attain population inversion.
The solution is to use a third metastable level, where the electrons can stay for longer duration. Under this
situation, the pumping will be between the other two levels and the electrons in the upper energy level will
quickly decay into the metastable level, leaving the upper level practically unpopulated at all times. The
transition from the metastable level to the ground level has a different frequency, which is the laser frequency.
The pumping frequency is between the upper level and the ground level. Thus the pumping is off-resonant to
the laser transition and it will not trigger the stimulated emission.
4.2.1.2 Three-Level System
Bloembergen proposed a mechanism where atoms are
pumped into an excited state by an external source of
energy, for example by an electric pulse or an optical
illumination. In addition to this excited state (say E3), the
system has a metastable state (say E2) and the atoms from
the upper level E3 decays spontaneously to this metastable
state and this transition is generally radiation less or non-
radiative (the energy being given away to the lattice). The
lifetime of the electrons in the metastable state E2 is such
that the rate of spontaneous decay from the upper level E3
to the ground level (say E1) is slower than the rate at which
the atoms decay from the upper level to the metastable
state, resulting in a population inversion between the
metastable level and the ground state (Fig. 4.5). The
population inversion can be achieved only by pumping into Fig. 4.5
a higher lying level, followed by a rapid radiative or non-radiative transfer into the upper laser level. This is
because in this way we can avoid the stimulated emission caused by the pump wave. The emitted photons
here are confined to a laser cavity to stimulate further the emission from the excited atoms. Larger width of
the excited level can make possible the absorption of a wider range of wavelengths to make pumping more
effective. Ruby laser works on the principle of a three-level system.
Since the lower level involved in the lasing (population
inversion) is the ground state of the atom, the three level
E4
system needs very high pumping power and yields low Fast decay
efficiency. Here more than half of the total number of the E3 Metastable state
Pumping
atoms have to be pumped to the excited state E3 before

achieving population inversion. The energy used to do this E Photon emission
in each of the cycle is wasted. However, the pumping power
Energy
E2
can be reduced significantly if the lower level involved
in the lasing is not the ground state. This will require at
Fa
least a four-level system (Fig. 4.6). Here, the pumping

s
td
ec
will transfer atoms from the ground state E1 to an excited

ay
state E4, from where they decay rapidly into the metastable
state E3 to make population N3 larger than population N2 to E1
achieve the condition of population inversion between E3
and E2 at moderate pumping. Population N
4.2.1.3 Four-Level System Fig. 4.6
The schematic of four-level system is depicted in Fig. 4.6 where four energy levels having energies E1,
E2, E3 and E4 with respective populations of N1, N2, N3 and N4 are shown. These energies follow the trend
E4 > E3 > E2 > E1. Here an optical pumping excites the atoms from the ground state E1 to the pump band
E4. The atoms from this level make a fast decay (radiationless transition) to the metastable energy level E3.
The population inversion of level E3 with the level E2 takes place when the lifetime of the transition from E3
to E2 is long compared to that of E4 to E3 (lasing level). The atoms in the metastable state E3 relax and start
to create laser transitions through spontaneous and stimulated emissions into energy level E2. The transition
from energy level E2 to the ground state (level E1) is fast just like level E4. This quickly de-excited atom leads
to a negligible population in the state E2 and maintains the population inversion. Since only a small number
of atoms need to be excited in the upper lasing level E3 to form the population inversion, a four-level laser
system is much more efficient and practical than the three-level laser system. The most popular four-level
solid state gain medium is Nd:YAG. All lasers based on neodymium-doped gain media are four-level lasers
except those operated on the ground state transition around 0.9–0.95 mm.
4.3 ChaRaCtERIStIC of LaSER LIght LO2

As discussed, laser radiation is achieved by the process of stimulated emission and the laser beam is highly
intense and directional. This radiation of a very pure frequency has the following main characteristics.
(i) Coherent: In simple words, the meaning of coherent is highly ordered. The word coherent comes
from another word “Cohero” which has the meaning “to stick together”. In fact, different parts
of the laser beam have a definite relationship to each other. This coherence is described in terms
of temporal coherence (coherence in time)
and spatial coherence (coherence in space)
(Fig. 4.7) which are required to produce high
quality interference.
Coherence in space
Ordinary light is not coherent because it comes
from independent atoms which emit on the time
scale of 10–8 seconds. A train of incoherent pho-
tons is shown in (Fig. 4.8) from which it is clear
that these photons are not in order, i.e., they do
not have a definite relationship with each other.
However, a degree of coherence can be found
in sources like the mercury green line, but their
coherence does not approach that of a laser.
Coherence in time
(ii) Monochromatic: The simple meaning of this Figure 4.7
word is that it is pure in colour or wavelength.
The light from a laser typically comes from one
atomic transition with a single precise wave-
length. So the laser light has a single spectral
colour and is almost the purest monochromatic
light available. It means the laser light is not
exactly monochromatic, but it has high de-
gree of monochromaticity. The deviation from
monochromaticity is due to the Doppler effect
of the moving atoms or molecules from which
the radiation originate.
(iii) Collimated: Collimated means it does not
spread out much. The light from a typical laser
emerges in an extremely thin beam with very Figure 4.8
little divergence, i.e, the beam is highly collimated. The high degree of collimation arises from the
fact that the cavity of the laser has very nearly parallel front and back mirrors as shown in (Fig. 4.9).
Because of this the light attains a parallel path after reflections from these mirrors. As it is clear from
the figure, the back mirror is made almost perfectly reflecting while the front mirror is about 99%
reflecting. Thus about 1% beam comes out from it, which we see as the output beam. Under this
process, however, the light passes back and forth between the mirrors many times in order to gain
intensity by the stimulated emission of more photons at the same wavelength. If the light is a bit off
axis, it will be lost from the beam.
Back Mirror Last Photon Front Mirror

Laser Light
Last Photon
100% Reflective 99% Reflective
Figure 4.9
The high degree of collimation or the directionality of a laser beam (single mode) is due to the
geometrical design of the laser cavity and to the fact that stimulated emission process produces
twin photons. A specific cavity design is shown in Fig. 4.10, where the angular spread of a beam
is signified by the angle q. In fact the cavity mirrors are shaped with concave surfaces towards the
cavity. This way the reflecting light is focused back into the cavity, which finally forms a beam waist
of radius r0 at one position in the cavity.
Wavefront Laser cavity
q
2q
Laser Beam
Fully Silvered Semi Silvered
Beam Waist Mirror
Mirror (radius r0) (Output
Mirror)
Figure 4.10
Considering the laser beam as the fundamental TEM00 mode (modes will be discussed in the chapter
on Electromagnetic Wave Propagation), the half angle beam spread can be, written as
l
q=
p r0
From this we obtain the angular spread as

0.637 l
2q =
r0
In addition to this, we can calculate the intensity, i.e., the power per unit area of a typical laser which is much
greater than other sources of electromagnetic radiation. This is due to the directionality and compactness of
the laser beam. In view of this, the intensity or irradiance of a laser beam in terms of its waist radius is given
by the following relation
P P
I= = 2 , where P is the power.
A p r0
4.4 MaIn CoMPonEntS of LaSER LO3
In order to understand the working principle of a laser, we should first know about the essential components
of the laser. These are given below
(i) Pumping: The method of raising the molecules or atoms from their lower energy state to higher
energy state is known as optical pumping. The optical pumping is needed for achieving population
inversion which is precondition for stimulated emission. In this case, the rate of stimulated emission
will exceed the rate of stimulated absorption. Hence, the intensity of light will increase during each
pass through the medium.
(ii) Active System: A system in which the population inversion is to be achieved is called as active
system or the gain medium for a laser. Laser systems are named based on the makeup of the gain
medium, which may be a gas, liquid or solid. The energy levels in the gain medium, those participate
in the radiation, determine the wavelength of laser radiation. Laser action has been observed in
over half of the known elements. Two of the most popular transitions in gases are 632.8 nm visible
radiation from neon and the 10.6 mm infrared radiation from the CO2 molecule.
(iii) Resonant Cavity: In a laser, the active system or the gain medium is enclosed in an optical cavity
(or resonant cavity) usually made up of two parallel surfaces, one of which is perfectly reflecting
reflector and the other surface is partially reflecting reflector. In this resonant cavity, the intensity of
photons is raised tremendously through stimulated emission process.
4.5 typEs Of LAsEr LO3
Nowadays different kinds of lasers are available, the most common being in a digital communications.
Virtually every house now has at least one – in their CD/DVD players and recorders. Some lasers can change
colour – they are called tunable lasers. The lasers now operate from the infrared to the ultraviolet regions.
Moreover, X-ray lasers are being developed using electron accelerators. The lasers now are available in the
wide range viz, solid lasers, liquid lasers, gas lasers, semiconductor lasers, etc.
4.5.1 Ruby Laser: Solid State Laser

Ruby laser is a solid state laser, which consists of three main parts (i) working material (ii) optical resonant
cavity and (iii) excitation source.
Working Material Ruby laser is made up of a crystal of ruby in the form of cylindrical rod having size 2
to 30 cm in length and 0.5 to 2 cm in diameter whose both ends are optically flat. One of the end is fully
silvered and other is partially silvered, so that they can act as fully and partially reflecting surfaces, respectively,
as shown in Fig. 4.11. Ruby rod is a crystal of Al2O3 Partially

in which chromium oxide is mixed as impurity so Xenon Flash Light Tube Silvered
Face (B)
that some of the Al3+ ions are replaced by Cr3+ ions. Fully
Silvered
These ‘impurity’ chromium ions give rise to the laser Face
action. Laser
(A)
Beam
The space between the two faces A and B is known
as the resonant cavity, in which the light (photon)
intensity can be built up by multiple reflections and
through stimulated emission. The ruby rod is wound Power Supply
by a helical xenon flash light tube with an excitation
source in the form of a power supply.
Inlet Outlet
Working Principle of Ruby Laser In this
laser, chromium ions are active centres which are Cooling Fluid (Liquid Nitrogen)
responsible for the laser transition. A simplified Figure 4.11
energy level diagram of chromium ions in ruby
crystal is shown in Fig. 4.12. In the normal state, most of the chromium ions are in the ground state E1. When
light from the flash tube of wavelength 5500 Å is made to fall upon the ruby rod, these incident photons are
absorbed by the chromium ions that rise to the excited state E3. Then they give a part of their energy to the
crystal structure and reach the metastable state, i.e., the E2 state. These ions in metastable state can remain for
a longer duration 10–3 sec. Therefore, the number of ions in this state goes on increasing while at the same
time number of ions in ground state goes on decreasing due to the optical pumping. Thus, the population
inversion is established between the metastable state and the ground state.
Short-live state
E3 10–8
Radiation-less Transition
Optical
Pumping
Metastable state
E2 10-3
6943 Å
5500 Å 6943 Å
Spontaneous
Emission 6943 Å
E1
Ground state
Figure 4.12
When an excited ion passes spontaneously from the metastable state to the ground state, it emits a photon of
wavelength 6943 Å. This photon travels parallel to the axis of ruby rod and stimulates the surrounding ions
present in the metastable state then by stimulated emission other photons are emitted, which are in the phase
with the stimulating photons. By successive reflections of these photons at the ends of the rod, every time
the stimulated emission is achieved, we obtain an intense, coherent and unidirectional laser beam from the
partially silvered face B.
The ruby laser operates at about 1% efficiency. It may produces a laser beam of 1 mm to 25 mm in diameter.
The beam obtained is in the form of pulses. However, on the advantage side, very strong beam as strong as
10,000 Watt in power is produced. Furthermore, the construction of this laser is simple and the operation is
very easy. For this reason, this laser is also known as practical laser. Other examples of solid state lasers are
Neodynium-YAG (Nd-YAG), Neodynium-Glass (Nd-Glass) and semiconductor lasers.
4.5.2 nd-Yag Laser: Solid State Laser

This laser is capable of producing very high power emissions, as a result of its lasing medium operates as four
level systems. The schematic of Nd-YAG laser is shown in Fig. 4.13. The lasing medium in the Nd-YAG laser
is colourless, isotropic crystal called Yttrium aluminium garnet (YAG-Y2Al5O12). The main dopant in the
lasing medium is Neodymium (Nd3+). When it is used in laser, Neodymium replaces 1% of Yttrium and the
crystal takes a light blue colour. The YAG has a relatively high thermal conductivity, which improves thermal
dissipation in thermal cavity. So continuous wave operation up to a few hundred watts is possible. Average
power of up to 1 kW is available when it is operated in pulse mode.
Fully Semi
Reflector Flash Lamp Reflector
Nd-YAG Crystal
Laser
Output
Flash Lamp
Power
Supply
Figure 4.13
The energy level diagram for Nd-YAG is shown in Fig. 4.14. These levels arise from three inner shell 4f
electrons of the Nd3+ ion, which are effectively screened by eight outer electrons (5S2 and 5P6). For the
operation of Nd-YAG lasers a cooling system is required. A Nd-YAG laser produces 30 times as much waste
heat as laser output with an efficiency of about 3%. The waste heat must be removed in order to ensure
proper laser operation by flooding the optical compartment with water. However optical distortion and image
[4S ,4F ]
[4F ,2H ]
11,502 cm–1 R2
4
F
0.73 µm
R1
0.8 µm
11,414 cm–1
1.06 µm
2,526 cm–1
4
I
n = 2×103 cm–1
4
I 2,001 cm–1
Figure 4.14
problem is created due to absorption of significant amount of flash lamp energy by water. This problem can
be overcome by flowing water over the outside of the optical cavity and by encasing the lasing rod and flash
lamp with transparent cooling jacket.
An advantage of Nd-YAG laser is that by using Q-switching, laser beam pulse frequency and shape can
be tailored where a shutter moves rapidly in and out of the path of the beam. In this manner beam output
is interrupted until a high level of population inversion and energy storage is achieved in the resonator.
If the optical cavity is switched from no reflection (low Q) to near total reflection (high Q), the cycle
can be optimised to build up the maximum population inversion before the pulse is generated. This
way, we get a beam pulse with high energy up to 1 J and a short pulse period down to 10 ns is obtained.
Applications
(i) Nd-YAG is used in material processing such as welding and drilling.
(ii) It is also used in photo disruption of transparent membrane of pathological origin, which can appear
in the interior chamber of eye or for iridectomy and in endoscopic applications.
(iii) It is used in range finders and target designators used in military context, which use Q-switched
lasers.
(iv) In scientific applications the Q-switched lasers with their second harmonic (l = 532 nm), third
harmonic (l = 355 nm) and fourth harmonic (l = 266 nm) are used.
4.5.3 helium-neon Laser: gas Laser

As we know that the output beam of the Quartz Tube Partially Silvered Mirror
ruby laser is not continuous. To overcome
this drawback, the gas filled laser was made
by A. Javan, W. Bennett and D. Herriott in Mixture of Laser
He & Ne Gases Beam
1961. It consists of a quartz tube having the
size about 1.5 cm in diameter and about 1
meter in length. The both ends of the tube Fully Silvered Mirror
R.F. Generator
are sealed by optically plane and parallel
mirrors, one of them being partially Figure 4.15
silvered (90% reflective) and the other one
is fully silvered (100% reflective).
In this laser system, a quartz tube is filled with a mixture of helium and neon gases in the ratio 10:1
respectively, at a pressure of about 0.1 mm of mercury (Fig. 4.15). This mixture acts as the active medium.
Helium is pumped upto the excited state of 20.61 eV by the electric discharge. The energy level diagram of
He-Ne laser is shown in Fig. 4.16.
Here, it can be seen that the excited level of He at 20.61 eV is very close to a level in Ne at 20.66 eV. It is so close
that upon collision of a He and a Ne atom, the energy can be transferred from the He to the Ne atoms. Thus, the
excited He atoms do not return to their ground state by spontaneously emitting photons rather they transfer their
energy to the Ne atoms through collisions. As mentioned, such an energy transfer can take place when the two
colliding atoms have identical states. Thus, the He atoms help achieving a population inversion in the Ne atoms.
An excited Ne atom passes spontaneously from the metastable state at 20.66 eV to the excited state at 18.70 eV by
Metastable State
20.66 eV
20.61 eV 6328 Å
6328 Å
6328 Å
18.70 eV
Spontaneous
Emission
Radiation-less
Transition
Ground state
He Ne
Figure 4.16
emitting a photon of wavelength 6328 Å. This photon travels through the gas mixture parallel to the axis
of the tube and stimulates the surrounding Ne atoms present in the metastable state. This way we get other
photons that are in the phase with the stimulating photons. These photons are reflected forth and back by the
silvered ends and the number of photons gets amplified through stimulated emission every time. Finally, a
portion of these intensified photons passes through the partially silvered end.
The He-Ne laser is the most common and inexpensive gas laser. Usually it is constructed to operate in the
red light at 6328 Å and in the infrared at 15,230 Å. According to Garmire, an unfocused 1 mW – He Ne laser
has a brightness equal to sunshine on a clear day (~ 0.1 W/cm2) and is just as dangerous to stare at directly.
4.5.4 Carbon Dioxide gas Laser

It is one of the earliest high power molecular gas laser that uses carbon dioxide gas molecule. This optical
device is capable of continuous output powers above 10 kW. It is also capable of extremely high power pulse
operation. It consists of discharge tube of size of about 2.5 cm in diameter and 5.0 cm in length. The both
ends of the tube are sealed by optically plane and parallel mirrors, one of them being semi-silvered and other
one is fully silvered. (Fig. 4.17).
Exhaust
Mixture of CO2, N2, He
CO2 N2 He
Power Supply
Fully Silvered Semi-Silvered
Mirror Mirror
Figure 4.17
The carbon dioxide gas laser mixture contain 15% CO2, 15% N2 and 70% He at a pressure of few mm of Hg.
This mixture is fed into the discharge tube through flow loop which is connected at one end of the discharge
tube. The dc excitations source is used that produces electric discharge. In starting nitrogen molecules are
allowed to enter in the discharge tube. They get excited by collision with electrons. Then excited nitrogen
molecules flow into the whole volume of resonant cavity and collide with the unexcited CO2 molecules
and transfer their energy to the desired laser level (Fig. 4.18). Nitrogen (N2) and helium (He) improve the
efficiency of the laser action, while oscillations take place between two vibrational levels of CO2. Nitrogen
helps producing a large populations in upper level and helium helps removing population from lower energy
level. Related energy levels of N2 and CO2 molecules are shown in Fig. 4.18. The radiated photons travel
back and forth between the end mirrors and get further amplified. It exhibits laser action at several infrared
frequencies but none in the visible. For example, it radiates light at 10.6 mm is far infrared region. It is one
of the most efficient lasers, capable of operating at more than 30% efficiency. Hence, this laser is suitable
for industrial applications both in terms of energy efficiency and high output beam; particularly it is used for
welding and cutting.
E n erg y Tr a n sfe r
CO2(001)
Laser
10.6 µm
CO2(100)
CO2(020)
CO2(010)
Ground State CO2(000)

N2 CO2
Figure 4.18
4.5.5 Semiconductor Laser
Semiconductor laser differs from the solid state and gas lasers in many aspects. It has remarkably small size,
exhibits high efficiency and can be operated at low temperature. When the current is passed through a p-n
junction diode in forward bias, holes move from p-region to n-region and the electrons move from n-region
to p-region. These electrons and holes are recombined in the junction region and emit photons due to the
transition of electrons from the conduction band to the valence band. This results in stimulated radiation
coming from a very narrow region near the junction. The action is intensified by increasing the current and
decreasing the junction thickness.
Semiconductor laser is made up of an active layer of gallium arsenide (GaAs) of thickness 0.2 microns. This is
sandwiched in between a n-type GaAsAl and p-type GaAsAl layer as shown in Fig. 4.19. The resonant cavity
is provided by polishing opposite faces of the GaAs crystal and the pumping occurs by passing electrical
current from an ordinary source (Power Supply). From this system GaAs semiconductor, laser beams of
wavelength ranging from 7000 Å to 30,000 Å can be produced.
GaAsAl
p-type
0.2 µm GaAs
Power
Supply
GaAsAl n-type
Figure 4.19
4.5.6 advantages and Disadvantages of Ruby laser, he-ne laser and

Semiconductor laser
The merits and demerits of solid laser (Ruby), gas laser (He–Ne) and semiconductor laser are given in
Table 4.1.
Table 4.1
Ruby laser He-Ne laser Semiconductor laser
Advantages
1 Easy to construct and operate Easy to construct and operate Easy in operation
2 Very strong and intense beam Continuous beam Long life, highly monochromatic,
upto a power of 10 kW tunable and continuous beam
3 Beam diameter as large as Exceptionally monochromatic Excellent efficiency with very high
25 mm beam with high operation operation duration (20,000 hrs.)
duration (10,000 hrs.)
Disadvantage
Its laser beam is only pulse It has got very low power It has got low power about 200
like and its operation duration about 0.5 – 5mW mW
is very less (few hrs.)
4.6 aPPLICatIonS of LaSERS LO4

Lasers have many applications in science, industry and medicine, some of which are listed below:
(i) Lasers have been used to measure long distances, so they are very useful in surveying and ranging.
For this purpose, a fast laser pulse is sent to a corner reflector at the point to be measured and the
time of reflection is measured to get the distance.
(ii) Lasers are electromagnetic waves of very high intensity and can be used to study the laws of
interaction of atoms and molecules.
(iii) Lasers are suitable for communication and they have significant advantages because they are more
nearly monochromatic. This allows the pulse shape to be maintained better over long distances. So
communication can be sent at higher rates without overlap of the pulses.
(iv) Laser beams are highly intense and are used for welding, cutting of materials, machining and drilling
holes, etc. Generally, carbon dioxide laser are used for such purposes, as it carries large power.
(v) Lasers are used most successfully in eye surgery, treatment of dental decay and skin diseases.
(vi) The laser beam is used in recording of intensity as well as in holography.
(vii) Laser is used in heat treatments for hardening.
(viii) Lasers are used as barcode scanners in library and in supermarket.
(ix) Laser is used in printers (Laser printers).
(x) Lasers are used in photodiode detection.
4.7 LaSER CooLIng LO4

So far you have learnt that the laser radiation is highly intense, highly coherent and highly monochromatic
light. It is amazing that the field of laser radiation can be used to cool down atoms to very low temperature,
for example up to 10–9 K. This can be understood based on an atom which is traveling toward a laser beam
and absorbs a photon from the laser. In this situation, the atom will be slowed by the fact that the photon has
momentum p = E/c = h/l, where E is the energy, c is the speed of light, h is the Planck’s constant and l is the
wavelength associated with the photon. If we assume that a number of sodium atoms are freely moving in a
vacuum chamber at 300 K (room temperature), i.e., the rms velocity of the sodium atom is about 570 m/s. Then
the momentum of the sodium atom can be reduced by the amount of the momentum of the photon, if a laser is
tuned just below one of the sodium d-lines (5890 Å and 5896 Å, about 2.1 eV) and the sodium atom absorbs
a laser photon when traveling toward the laser. It would take a large number of such absorptions to cool the
sodium atoms to nearly 0 K. The change in speed from the absorption of one photon can be calculated from
Dp pphoton Dv pphoton
= = fi Dv =
p mv v m
The above expression shows a lot of photons, but according to Chu a laser can induce on the order of 107
absorptions per second so that an atom could be stopped in a matter of milliseconds.
There is a conceptual problem that an absorption can also speed up an atom if it catches it from behind. So
more absorptions from head-on photons are necessary to have. This is accomplished in practice by tuning the
laser slightly below the resonance absorption of a stationary sodium atom. More precisely, with the opposing
laser beams with perpendicular linear polarisations, atoms can be selectively driven or “optically pumped”
into the lower energy levels. This method of cooling sodium atoms was proposed by Theodore Hansch and
Arthur Schawlow at Stanford University in 1975 and was achieved by Chu at AT & T Bell Labs in 1985. Here
sodium atoms were cooled from a thermal beam at 500K to about 240 mK.
4.8 hoLogRaPhY LO5

Holography is one of the remarkable achievements of modern science and technology, which has been possible
only because of the lasers. The word “holography” was originated from the Greek words “holos” and “grapho”.
The meaning of “holos” is “whole” and of “grapho” is “write”. So holography means complete record of the
image. Holography is a three-dimensional (3D) laser photography. It is lensless photography in which an
image is captured as an interference pattern. The image thus obtained is called a hologram, which is true 3D
record of the object. Holography not only records the amplitude but also the phase of the light wave with the
help of interferometric techniques. This recorded reference pattern contains more information than a focused
image and enables the viewer to view a true 3D image, which exhibits parallax. The technique of holography
was invented by Gabor in 1947.
4.8.1 Principle of holography
In holography, there are two basic waves that come together to create the interference pattern. One wave is
called object wave and another wave is called reference wave. When an object wave meets a reference wave,
it creates a standing wave pattern of interference. This is then photographed, which we call a hologram.
4.8.2 Requirements of holography
Following are some requirements for the absolute holography.
(i) Since holography is an interference phenomenon, there should not be a path difference between the
object wave and the reference wave more than the coherence length. This is necessary to achieve
stable interference fringes.
(ii) Spatial coherence is important so that the reference wave and the scattered object waves from
different regions can interfere properly.
(iii) Since reconstructed image coordinates depend on wavelength as well as position of the reconstructing
source, it is necessary that the source emits a narrow band of wavelength and it is not broad in the
interest of obtaining good resolution in the reconstructed image.
(iv) In order to obtain aberrations free reconstructed image, it is necessary that the reconstructing source
is of the same wavelength and is situated at the same position with respect to the hologram as the
reference source.
(v) All recording arrangement like film, object, mirrors etc., must be motionless during the exposure.
4.9 hoLogRaPhY VERSuS ConVEntIonaL PhotogRaPhY LO5

Holography represents a photographic process in a broad sense, but essentially it differs from a usual photo, as the
phase of light waves scattered by the object carries the complete information about 3D structure of the object. A
conventional photography is a 2D image of a 3D scene, which brings into focus every part of the scene that falls
within the depth of the field of the lens. Due to this, a conventional photograph lacks the perception of the depth
or the parallax with which we view a real life scene. Since a conventional photograph only records the intensity
pattern, 3D character of the object scene is lost. Contrary to this, the hologram contains depth and parallax, which
provides the ability to see around the object to objects placed behind. It gives information about amplitude as well
as the phase of an object. So hologram preserves information about the object for latter observation.
In conventional photography, there is one to one relationship between object and image point as the light
originating from a particular point of scene is collected by a lens focused on that particular point. However,
in holography lens is not used and this is a complex interference pattern of microscopically spaced fringes.
Hologram receives light from every point of a scene and hence there is no one to one relationship. This is a
record of entire signal wave.
In conventional photography, radiated energy is recorded and phase relationship of wave arriving from
different distances and directions is lost. However, in holography phase relationship is recorded by using the
technique of interference of light waves.
4.10 RECoRDIng anD REConStRuCtIon of IMagE on hoLogRaPh LO5

Recording of a hologram is a result of superpositions of the object wave and the reference wave, which is
usually a plane wave. This interference pattern is recorded by a photographic plate that contains information
about amplitude as well as phase of the object wave. In order to see the image, hologram is illuminated with
another wave called the reconstruction wave which is identical to the reference wave in most of the cases.
This is called reconstruction of image on the hologram.
4.10.1 theory
If the object is a point scatterer and it is made of large number of such points, then the composite wave
reflected by the object will be the vectorial sum of all object waves scattered from are these points. As
mentioned earlier, holography records the object wave, particularly the phase (say j) associated with it. So
we can represent the object wave, which is due to the superposition of waves from point scatterers on the
object, as
Y1(x, z) = A1(x, z) cos (f – wt) (i)
where w is the frequency. The object wave represented by Eq. (i) lies in the plane of photographic plate at
y = 0.
Now, we consider a reference wave propagating in the xy plane and inclined at an angle a from the y axis. In
view of this, the field associated with the reference wave can be written as
 
Y2(x, y, z) = A2 cos (k ◊ r - w t )
= A2 cos (kx sin a + ky cos a – wt) (ii)
At the photographic plate, i.e., at y = 0, this field becomes
Y2(x, z) = A2 cos (kx sin a – wt).
sin a
Since the propagation constant k = 2p/l, kx sin a = 2p x
l
Here sin a/l is defined as the spatial frequency (say x). So the field associated with the reference wave
becomes
Y2 (x, z) = A2 cos (2px x – wt) (iii)
A comparison of equation (iii) with equation (i) yields that the phase linearly varies with x.
Simple method of superposition enables us to calculate the total field at the photographic plate (at y = 0) as
Y = Y1 + Y2
Y(x, z, t) = A1(x, z) cos (f – wt) + A2 cos (2pxx – wt) (iv)
In view of the response of photographic plate to the intensity we find below the measure of intensity pattern
recorded by the photographic plate as
I(x, z) = Average value of Y2(x, z, t)
= < Y2 (x, z, t) >
or I(x, z) = A12 (x, z) < cos2 (f – wt)> + A22 < cos2 (2pxx – wt)
+ 2A1(x, z) A2 <cos (f – wt) cos (2pxx – wt> (v)
2
As we know that < cos (f – wt) > = ½,
<cos2 (2pxx – wt) > = ½,
< 2cos (f – wt) cos (2pxx – wt) >
= ½ <cos (f + 2pxx – 2wt) + cos (f – 2pxx) >
[Using 2 cos q1 cos q2 = cos (q1 + q2) + cos (q1 – q2)]
The average value of cos (f + 2pxx – 2wt) can be obtained by using simple integration
1 2p / w
2p Ú0 cos (f + 2px x - 2w t ) dt ,
as the average value of cos wt over the period T = 2p/w

1 2p / w
=
2p Ú0 cos w t dt.
So it comes out to be zero. With this, the intensity I(x, z) is written as

I(x, y) = A12 (x, z)/2 + A22 /2 + A1(x, z) A2 cos (f – 2pxx) (vi)
The above equation shows that the intensity I is the function of phase f(x, z). It means the phase information
of the object wave is recorded in the intensity pattern.
In order to obtain a hologram, we develop the photographic plate containing above intensity pattern. In this
context, the ratio of the transmitted field to the incident field is defined as the transmittance of the hologram
that depends on I(x, z). Using a suitable developing process, the condition under which the transmittance is
linearly related to I(x, z) can be obtained. Under this condition, if Re(x, z) denotes the field of the reconstruction
wave, at the hologram plane, then the transmitted field can be taken as
Te(x, z) a Re(x, z) I(x, z).
Taking Kp as the proportionality coefficient and putting the value of I(x, z) from equation (vi), we obtain
È 2 2 ˘
Te ( x, z ) = K p Re ( x, z ) Í A1 ( x, z ) + A2 + A1 ( x, z ) A2 cos (f - 2px x) ˙ (vii)
Î 2 2 ˚
In the case when the reconstruction wave is identical to the reference wave Y2(x, z), the above equation
becomes
È 2 2˘
Te ( x, z ) = K p A2 Í A1 ( x, z ) + A2 ˙ cos (2px x - w t )
Î 2 2 ˚
+ K p A22 A1 ( x, z ) ◊ cos (2px x - w t ) cos (f - 2px x)
Again using 2 cos q1 cos q2 = cos (q1 + q2) + cos (q1 – q2), we get the following expression for Te(x, z)
È 2 2˘
Te ( x, z ) = K p A2 Í A1 ( x, z ) + A2 ˙ cos (2px x - w t )
Î 2 2 ˚
K p A22 A1 ( x, z )
+ cos (4px x - f - w t ) ◊ cos(f - w t ) (viii)
2
The above equation contains three terms, which may be analysed as follows.
(i) First term being proportional to A22 represents the reconstruction wave whose amplitude is modulated
by the term A12 (x, z), i.e., by the amplitude of object wave. The factor cos (2pxx – wt) shows that this
part of the total field is traveling in the direction of the reference wave.
(ii) The second term is identical to equation (i) within a constant term. Hence, this represents the original
object wave. Having appeared in transmitted field, it gives rise to a virtual image.
(iii) The third term carries the phase f(x, z) in addition to the term 4pxx, but with negative sign. It means
this wave has a curvature opposite to the object wave, i.e., if the object wave is diverging spherical
wave, then the last term (third term) shows a converging spherical wave. Hence, this wave forms a
real image of the object contrary to the second term. This image can be photographed with the help
of a film.
4.11 tYPES of hoLogRaMS LO6

In order to construct a hologram, we need two coherent light waves, one is the object wave carrying
information about the object and the other is a plane wave that is called reference wave. There are various
types of holograms, but the most common ones are the transmission hologram and the reflection holograms.
4.11.1 transmission hologram
This type of hologram is commonly used. If the object wave and the reference wave emerge from the same side
of the holographic film, then the hologram is called transmission hologram (Fig. 4.20). Another characteristic
of transmission hologram is the low diffraction efficiency and weak image reconstruction.
Laser
Object
Beam
Mirror
Beam
Splitter
Film
Object
Reference
Beam
Figure 4.20
4.11.1.1 Recording Process
As mentioned, in order to make a hologram, two coherent light waves (laser light) are required (Fig. 4.21).
The first one is called the object wave, which is reflected from the object and carries information about the
object. The second one is called reference wave and is a plane wave without information. These two waves
generate an interference pattern, which is recorded in the form of a hologram on film emulsion. For obtaining
the stable interference patterns, absolutely stable conditions are required during the exposure of the film. This
recorded hologram is called transmission hologram because the light passes through the holographic plate.
Laser
Laser
Spatial filter Spatial filter

Shutter
Beam splitter Beam splitter Mirror
Object
Spatial filter Virtual image
Hologram
Filmplate
Mirror Eyes
4.11.1.2 Reconstruction Process
We can reconstruct the holographic image by developing the hologram and then placing it in its original
position in the reference beam as during its recording. If we look along the reconstructed object wave, we see
a replica of the object and as we shift viewpoints we see object from different perspectives. Thus, the object
appears to be three-dimensional. During the reconstruction of the transmission hologram, the light does not
pass through the image, but it creates a wavefront that makes it appear as though the light had been generated
in the position of the object. This image thus formed is called virtual image (Fig. 4.22). Contrary to this, an
image having light actually passing through it is called a real image.
4.11.1.3 Properties
Some important properties of transmission hologram are as follows.
(i) When viewed with white light, the transmission holograms look like a blurry rainbow image.
(ii) These holograms are viewed as sharp images when we use the shining laser light through the hologram.
(iii) Less resolving power is needed in materials.
(iv) Transmission hologram can be formed in a simple setup.
(v) Greater depth of the scene is possible in transmission holograms.
4.11.2 Reflection hologram
The holograms that are viewed with white light source on the same side as the viewer are known as reflection
holograms. In such a hologram, a truly three-dimensional image is seen near its surface. This hologram is the
most common type shown in galleries. The light is located on the viewer’s side of the hologram at a specific
angle and distance. The image thus formed consists of light reflected by the hologram. There are two types
of reflection holograms.
Recording reflection hologram
4.11.2.1 One-Step Hologram
Laser
Here, the resolution of film emulsion is high, as the recording
of reflection hologram needs 10 to 100 times much power than a Shutter
transmission hologram. Thus, exposure time is long. During the Beam splitter Mirror
Object
process of recording the hologram, the two waves namely the reference Spatial filter
wave and the object wave illuminate the film plate on opposite sides
(Fig. 4.23). In this case, the fringes are formed in layers and are more
or less parallel to the surface of the emulsion. If a highly directed Filmplate
Mirror
beam of white light illuminates a reflection hologram, it selects the
appropriate band of wavelengths to reconstruct the image and the
Spatial filter
remainder of the light passes straight through.
Mirror
4.11.2.2 Two-Step Hologram
Figure 4.23
This hologram involves two steps. First we make a transmission
hologram called H1 (Fig. 4.24). This is called a master or first
hologram. We make multiple copies from the master hologram. We
make transfer copies of master hologram. Transfer copy means making Laser
another hologram using the image on the master as the subject. These
transfer holograms are either laser-visible transmission holograms or Shutter
Mirror
reflection holograms H2. Suppose we want any object in the final Beam splitter
hologram just to appear half in front and half behind the recording Spatial filter Spatial filter
plate. In such circumstances, the two step hologram is of great use.

4.11.2.3 Properties
Some important properties of reflection hologram are as follows.
(i) These holograms can be viewed in regular light. Mirror Mirror
H1 H2
(ii) The finished reflection hologram is monochromatic.
Figure 4.24
4.11.3 White Light hologram: Rainbow hologram
Rainbow holograms that can be viewed in white light and produce 3D images are very popular holograms. A
double holographic process makes them, in which an ordinary hologram is used as the object and a second
hologram is made through a slit. A horizontal slit limits the vertical perspective of the first image so that
there is no vertical parallax. The coherence requirement can be removed by using slit process. So the image
brightness is obtained from ordinary room light while maintaining the 3D character of the image as the
viewer eye is moved horizontally. If viewer eye is moved vertically, no parallax is seen and the image colour
sweeps through the rainbow spectrum from blue to red.
4.11.3.1 Embossed Hologram
Many variations of hologram can be made between the reflection and transmission types of holograms.
Embossed hologram is one of such types of holograms, which is used widely in most security applications.
These holograms offer an effective method of protection against any forms of manipulation, as they are too
difficult to copy due to their complex structure. All credit cards and passports have embossed hologram.
In this hologram, the original hologram is recorded in a photosensitive material called photoresist. These
holograms are easily produced at large scale and also at a very low cost.
4.11.3.2 Volume Hologram
Volume holograms are produced when the thickness of the recording material is much larger than the light
wavelength used for recording. These are transmission holograms and are also known as thick holograms,
which are mainly considered as a high-density data storage technology. These are 3D holograms created by
recording the interference pattern of two mutually coherent light waves. The angle of difference between the
object wave and the reference wave is 90° to 180°. Due to certain unique properties, volume holograms are
used widely in various spectroscopic and imaging applications.
4.12 aPPLICatIonS of hoLogRaPhY LO7

Holography represents examples of recombining of scattered radiation. It is a product of interference of light,
which is used to measure very small optical path lengths with precision by using wavelength of light and
interference. Now holography is being used in industry, communication and other engineering problems also.
You would have seen hologram on tickets, original covers of software programs, credit cards etc. This is used
to prevent falsification. Another important application is through bar code readers used in shops, warehouses,
libraries and so on. In aircraft industry, holography technology is used through head up displays (HUD) which
help the pilot to see instrument panel on to the windscreen.
Some other important applications of holography are given below.
4.12.1 time average holographic Interferometry

This interferometry is very useful for determining or studying the modes of vibration of complex structures.
Hologram is prepared using a long exposure time than the periods of vibrations being studied. This hologram
freezes many images, mapping the motion of vibrating surface. Interference fringes pattern provides
information about the relative vibrational amplitudes as a function of position on the surface.
4.12.2 Microscopy
A hologram contains many separate observations of microscopic particles. Image provided by hologram may
be viewed by focusing on any depth of unchanging field. Microscopic hologram is made by illuminating the
specimen by laser light, a part of which is split off outside the microscope and is routed to the photographic
plate to rejoin the subject beam processed by the microscope. It can be shown that if lr > ls, where lr is the
wavelength of reconstructing light and ls is the wavelength used in holography, then the magnification is
Ê v ˆ Ê lr ˆ
M =Á ˜Á ˜
Ë u ¯ Ë ls ¯
Here u is the object distance from the film and v is the corresponding image distance from the hologram.
However, these distances are equal, i.e., u = v, if the reference and reconstructed wavefronts are both plane
wave.
4.12.3 ultrasonic hologram
As the words “ultrasonic holograms” suggest, the waves producing a hologram may not necessarily be
electromagnetic in nature. Also, the holographic principles do not depend on the transverse nature of the
radiation. Holograms generated with the help of ultrasonic waves are very useful because of the ability of
such waves to penetrate the objects that are opaque to visible light. Holograms formed by ultrasonic waves
are very useful to get 3D images inside the opaque bodies.
4.12.4 holocameras
Hologram can be developed and viewed with the help of holocameras, which do not use photographic film.
Thermoplastic recording material is used in holocameras and image development is done by electrical and
thermal means. The image development does not need wet chemical processing. Also, it can be completed in
a few seconds without repositioning the recording.
4.12.5 holographic Data Storage

Data can be stored by holographic technique. It is very interesting that the data can be reduced to dimensions
of the order of wavelength of light. Therefore, volume holograms can be useful to record vast quantities of
information. Photosensitive crystal like potassium bromide with colour centres or lithium niobate are used in
place of thick layered photoemulsion. Small rotation of crystal takes place of turning pages.
s UmmarY
The main topics discussed in this chapter are summarized below.
✦ Laser was introduced as a special type of device that amplifies light and produces a highly intense and
highly directional beam which mostly has a very pure frequency.
✦ It was made clear the population inversion is the basic requirement for the operation of the laser.
✦ For achieving the laser radiation, the concept of stimulated emission was discussed in detail along with
the inclusion of Einstein’s coefficients.
✦ The main components of laser were discussed and based on the gain medium the lasers were classified
as solid state laser, gas laser or semiconductor laser.
✦ Ruby laser, Nd-YAG laser, He-Ne laser, CO2 laser and semiconductor laser were discussed in detail and
the energy diagrams provided.
✦ It was mentioned that the lasers have diverse applications in different fields of science and technology.
These applications were talked about in brief.
✦ A new concept of laser cooling was discussed in detail. It was shown how a highly intense and coherent
light of laser can cool the sodium atoms to 10–6 K.
✦ Another exciting filed of holography was introduced and it was mentioned that with the help of lasers
the holograms can be developed that give 3D picture of the objects.
✦ Principle and the requirements of the holography were discussed.
✦ The advance/additional features of holography from those of conventional photography were talked
about.
✦ Detailed description of recording and reconstruction of image on holograph were discussed.
✦ Two types of holograms, namely transmission holograms and reflection holograms, were discussed in
detail along with their recording and reconstruction processes and the properties.
✦ White light hologram was introduced, which is also known as rainbow hologram. Then the embossed
and volume holograms were talked about.
✦ Various applications of holography were discussed including time average holographic interferometry,
microscopy, ultrasonic holograms, holocameras and the holographic data storage.
s olVeD e XamPles
E xamplE 1 Determine the energy and momentum of a photon of a laser beam of wavelength 6328Å (Given:
h = 6.63 ¥ 10–34 J K sec. and c = 3.0 ¥ 108 m/sec).
Solution Given l = 6328 ¥ 10–10 m, h = 6.63 ¥ 10–34 J K sec. and c = 3 ¥ 108 m/sec.
hc
Formula used E = hn =
l
6.63 ¥ 10-34 ¥ 3 ¥ 108
= = 1.05 ¥ 10-19 Joule
6.328 ¥ 10-7 m
E = 3.143 Joule
E h 6.63 ¥ 10-34
Momentum p = = = = 1.05 ¥ 10-27 kg ◊ m/sec
c l 6.328 ¥ 10-7
p = 1.05 kg m/sec.
E xamplE 2 Calculate the energy of laser pulse in a ruby laser for 2.8 ¥ 1019 Cr3+ ions. If the laser emits
radiation of wavelength 6943Å.
Solution Given: l = 6943 ¥ 10–10 m, n = 2.8 ¥ 1019
The energy of a photon, = hn
and the total energy due to n Cr3+ ions is
hc 6.63 ¥ 10-34 ¥ 3 ¥ 108
E = nhn = n = 2.8 ¥ 1019 ◊
l 6.943 ¥ 10-7
8.02 J
E xamplE 3 A three-level laser emits a light of wavelength of 5500 Å, What will be the ratio of population
of upper level (E2) to the lower energy level (E1) if the optical pumping mechanism is shut off (Assume
T = 300 K).
At what temperature for the conditions of (a) would the ratio of populations be 1/2?
Solution Given l = 5500 Å
Formula used is
hc
E2 - E1 = hv =
l
(6.63 ¥ 10-34 J/sec) ¥ (3 ¥ 108 m/sec)
=
(5.5 ¥ 10-7 m) ¥ (1.6 ¥ 10-19 J/eV)
= 2.26 eV
and kT can be calculated as

kT = (8.62 ¥ 10–5 eV/K) ¥ (300 K)
= 0.0259 eV
The ratio of upper to the lower energy levels i.e.,
E2
= e( - E2 - E1 )/ kT = e - 2.26/0.0259
E1
= e -87.3
E2
= 1.3 ¥ 10 –38
E1
Given
1
E2 / E1 =
2
Then by using above equation
E2 1
= = e - (E2 - E1)/kT
E1 2
E2 - E1
or e(E2 - E1)/ kT = 2 or = log e2
kT
E2 - E1 2.26 eV
or T= = = 37832.75 K
K log e 2 Ê -5 eV ˆ
8.62 ¥ 10 ¥ (0.693)
ËÁ K¯
˜
or T = 37832 K
This temperature is much hotter than the sun.
E xamplE 4 (a) A He-Ne laser of wavelength 6328 Å has an internal beam of radius 0.23 mm. What would
be the beam divergence angle?
(b) What lower limit might be expected for the beam divergence, if we can control the beam-
waist radius (r0) by lower cavity design and selecting the wavelength. By what factor will
the beam divergence decrease if we design a laser having a beam waist of 2.4 mm radius and
wavelength 2000 Å?
Solution
(a) Given l = 6328 Å, r = 0.23 mm
Formula used is
l 6328 ¥ 10-10 m
q= = = 8.76 ¥ 10 –4 rad.
p r0 3.14 ¥ 2.3 ¥ 10-4 m
i.e., beam radius increases about 8.76 cm at every 100 m distance.
(b) Given, l = 2000 Å, r = 2.4 mm
then
l 2000 ¥ 10-10 m
q= = = 2.65 ¥ 10 –5 rad.
p r0 3.14 ¥ 2.4 ¥ 10-3 m
i.e., about 33 fold decrease in beam spread over the He-Ne laser described in part (a) and in this case beam radius
increases about 2.65 mm at every 100 m distance.
E xamplE 5 A pulsed ruby laser consists of ruby crystal in the form of a cylinder of size 6.0 cm in length
and 1.0 cm in diameter. Ruby laser is made of Al2O3 crystal in our case one aluminium ion in every 3500 has
been replaced by chromium ion Cr3+ ion and these same ions also produce laser light which occurs by three
level mechanism at a wavelength of 6944 Å. [Given density (r) of Al2O3 = 3700 kg/m3 and Molar mass =
0.102 kg/mol.]
Solution Given, length (l) = 6.0 ¥ 10–2 m, diameter (D) = 1.0 ¥ 10–2 m, l = 6944 Å, density (r) of Al2O3 = 3700kg/m3,
Molar Mass M = 0.102 kg/Mol.
Formula used for no. of aluminium ions is
2 N.m 2N ◊ r ◊ V
N Al = =
M M
where m is the mass of ruby cylinder and factor 2 accounts for two aluminium ions in each molecule of Al2O3. The
volume V is given as
2
Ê Dˆ p
V = p r 2l = p Á ˜ l = D 2l
Ë 2¯ 4
1
= ¥ 3.14 ¥ (1.0 ¥ 10-2 ) 2 ¥ 6.0 ¥ 10-2
4
= 4.7 ¥ 10- 6 m3
Thus,
2 ¥ (6.0 ¥ 1023 per mol) ¥ (3.7 ¥ 103 kg/m3 ) ¥ 4.7 ¥ 10-6 m3
N Al =
0.102
23
= 2.1 ¥ 10
and the number of chromium ions Cr3+ ions is given by
N Al
N cr = = 6.0 ¥ 1019
3500
The energy of the stimulated emission photon is given by
hc 4.1 ¥ 10-15 eV. sec ¥ 3 ¥ 108 m/sec

E = hv = =
l 6.944 ¥ 10-7 m
= 1.8 eV.
And hence the total energy due to all the pulses is given by
Etotal = Ncr ◊ E
= 6.0 ¥ 1019 ¥ 1.8 eV ¥ 1.6 ¥ 10–19 J/eV = 17 Joules
E xamplE 6 Calculate the power per unit area delivered by a laser pulse of energy 4.0 ¥ 10–3 Joule, the pulse
length in time as 10–9 sec and when the pulse is focused on target to a very small spot of radius 1.5 ¥ 10–5 m.
Solution Given P = 4.0 ¥ 10–3 J, r = 1.5 ¥ 10–5 m
Formula used for power delivered per unit area is given by
P 4.0 ¥ 10-3 J
I= , where P =
A 10-9 sec.
or P = 4.0 ¥ 106 W
and A = pr2 = 3.14 ¥ (1.5 ¥ 10–5)2 = 7.065 ¥ 10–10 m2

P 4.0 ¥ 106 W
so I= = = 5.7 ¥ 1015 W/m 2
A 7.065 ¥ 10-10 m 2
or I = 5.7 ¥ 1015 W/m2
E xamplE 7 A laser beam has wavelength of 7200 Å and aperture 5 ¥ 10–3. The laser beam is sent to moon
at a distance 4 ¥ 108 m from the earth. Determine (a) angular spread and (b) a real spread when it reaches
the moon.
Solution Given l = 7.2 ¥ 10–7 m,
d
radius r = = 2.5 ¥ 10–3 m, D = 4.0 ¥ 108 m
2
Formula used is
(a) Angular spread (q)
0.637 l
=
r
0.637 ¥ 7.2 ¥ 10-7
q=
2.5 ¥ 10-3
or = 1.834 ¥ 10- 4 radian
(b) Areal spread = (qD)2 = (4 ¥ 108 ¥ 1.834 ¥ 10–4)2

= 53.85 ¥ 108 m2
E xamplE 8 A 0.1 W laser beam with an aperture of 5.0 mm emits a light of wavelength 6943 Å. Calculate
the areal spread and intensity of the image when the beam is focused with a lens having focal length 100 mm.
Solution Given:
diameter
radius of aperture =
2
or r = 2.5 ¥ 10 m, l = 6.943 ¥ 10–7 m, f = 0.1 m, P = 0.1 W
–3
Formula used is
0.637 l
Angular spread (q ) =
r
0.637 ¥ 6.943 ¥ 10-7 m
q=
2.5 ¥ 10-3
or q = 1.769 ¥ 10–4 radius
A real spread = (q ◊ D)2 = (q ◊ f)2 (\ D = f)
= (1.769 ¥ 10–4 ¥ 0.1 mm)2
= 3.129 ¥ 10–10 m2
and the intensity is given by
Power ( P ) P 0.1 W
I= = =
Area ( A) A 3.129 ¥ 10-10 m 2
= 3.196 ¥ 108 W /m 2
E xamplE 9 For an ordinary source, the coherence time tc = 10–10 sec. Obtain the degree of non-monochro-
maticity for lo = 5400 Å.
Solution Given tc = 10–10 sec
1 1
Dv = = = 1010 Hz
t c 10-10
c 3.0 ¥ 108 1
For l0 = 5400 Å, n 0 = = = ¥ 1016
l0 5400 ¥ 10-19 18
degree of non-monochromaticity
Dn 18 ¥ 1010
= = 18 ¥ 10-6 = 0.000018
n0 1016
o bJectiVe t YPe Q UestioNs
Q.1 LASER is a short form of

(a) Light Amplification Stimulated Emission Radiation
(b) Light Amplification by Stimulated Emission of Radiation
(c) Light Absorption by Stimulated Emission of Radiation
(d) Light Absorption by Spontaneous Emission of Radiation
Q.2 Mention the process under which an electron jumps from higher energy state to lower energy state by
the influence of incident photon
(a) induced emission (b) spontaneous emission
(c) simple emission (d) none of these.
Q.3 Laser beam is
(a) highly monochromatic (b) highly coherent
(c) highly collimated (d) all of these.
Q.4 What is the life-time of electron in metastable state?
(a) 10–3 sec (b) 10–5 sec (c) 10–8 sec (d) 10–7 sec
Q.5 The number of atoms in the higher energy state is larger than lower energy state. This state is known as
(a) metastable state (b) ordinary state (c) excited state (c) none of these.
Q.6 In the population inversion
(a) the number of electrons in higher energy state is more than the ground state
(b) the number of electrons in lower energy state is more than higher energy state
(c) the number of electrons in higher and lower energy state are same
(d) none of them.
Q.7 The relations between Einstein’s coefficient A and B is
3
8p hn 3 8p 2 h 2n 3 Ê 2p hn ˆ 8p hc
(a) (b) (c) Á (d)
c 3
c3 Ë c ˜¯ l
Q.8 Laser beam is made of
(a) electrons (b) highly coherent photons
(c) very light and elastic particles (d) none of them
Q.9 In ruby laser which ions give rise to the laser action?
(a) Al2O3 (b) Al3+ (c) Cr3+ (d) none of them
Q.10 The output beam in ruby laser is
(a) continuous (b) discontinuous (c) both (a) & (b) (d) none of these
Q.11 Which one of the following laser have highest efficiency, ruby, He-Ne and semiconductor and carbon
dioxide?
(a) ruby (b) semiconductor (c) He–Ne (d) carbon-dioxide
Q.12 The He–Ne laser produces the laser beam of wavelengths
(a) 6943Å (b) 6328Å (c) 6320Å (d) 6940Å.
Q.13 In He–Ne laser the ratio of the He to Ne is
(a) 10:1 (b) 1:10 (c) 100:1 (d) none of these.
Q.14 The method of population inversion to the laser action in He–Ne laser is:
(a) molecular collision (b) direction conversion
(c) electric discharge (d) electron impact.
Q.15 Ruby laser produces the laser beam of wavelength
(a) 6943 Å (b) 6328 Å (c) 6320 Å (d) 6940 Å.
Q.16 Characteristics of laser beam are
(a) highly directional (b) highly intense
(c) highly monochromatic (d) all of them.
Q.17 Holography was discovered by Dennis Gabor in
(a) 1948 (b) 1847 (c) 1748 (d) none of these.
Q.18 Holography records intensities and phases of light coming from an object on holographic plate has
(a) complete information of object (b) incomplete information of object
(c) no information of object (d) none of these.
Q.19 Holography produces the image
(a) real (b) virtual (c) both (a) & (b) (d) none of these.
Q.20 Which of the following statement is correct?
(a) Holography has been used to see the working condition of inner organs of the body in three
dimension
(b) data storage
(c) in non-destructives testing of materials
(d) all of these.
Q.21 Information carrying capacity of hologram is
(a) large (b) small (c) zero (d) none of these.
P ractice P roblems
Q.1 What do you mean by laser and its working principle, important requirements and applications?
Q.2 (a) Explain the term ‘absorption’, ‘spontaneous’ and ‘stimulated’ emission of radiation. Obtain a
relation between transition probabilities of spontaneous and stimulated emission.
(b) What are Einstein’s coefficient? Derive Einstein relation.
Q.3 Explain the construction and working principle of Ruby laser.
Q.4 What is the principle of laser?

Q.5 Discuss salient characteristics of laser beam.
Q.6 Describe various applications of laser?
Q.7 Explain the characteristics of laser beam. What are the necessary conditions for Lasing action?
Q.8 Explain the concept of directionality and monochromaticity as applied to lasers.
Q.9 Discuss Einstein’s coefficients. Derive relation between them.
Q.10 Explain the terms: spontaneous and stimulated emission, population inversion, optical pumping.
Q.11 What is meant by population inversion and how is it achieved in practice.
Q.12 (i) What do you mean by spontaneous emission? Explain in detail various properties of Laser which
makes it a superior source to study the phenomenon of interference.
(ii) Explain the working of Carbon-dioxide laser.
(iii) What is Nd YAG laser and explain its working principle?
Q.13 Define the following terms.
(a) Population inversion (b) Pumping (c) Active system
Q.14 What do you understand by solid state laser? Describe the principle, construction and working of Ruby
Laser.
Q.15 Explain with the help of a neat diagram the working of Ruby laser.
Q.16 Explain the concept of coherence in lasers. What are necessary conditions for Lasing action? Give
main components of a Ruby or He-Ne laser and principle of laser action.
Q.17 Explain the terms stimulated emission of radiation and optical pumping. Explain how lasers can be
produced by He–Ne gas. How is it superior to a Ruby laser?
Q.18 Explain with the help of a neat diagram the principle and working of a He-Ne laser.
Q.19 What are the specialties of laser light? Give the description of He–Ne laser.
Q.20 Write a note on He–Ne laser.
Q.21 (a) Discuss with suitable diagrams the principle, construction and working of Helium–Neon laser.
(b) Describe the various applications of lasers.
Q.22 Discuss the essential requirements for producing laser action. Describe a He–Ne laser.
Q.23 What conditions must be fulfilled for a semiconductor laser? Explain.
Q.24 Give description of semiconductor laser and discuss the characteristics of laser beam.
Q.25 Describe the construction and working of a semiconductor laser with necessary diagrams.
Q.26 How does a semiconductor laser differ from other lasers? Explain main features of the semiconductor
laser.
Q.27 Mention how does a semiconductor laser differ from Ruby laser.
Q.28 Write a note on the following.
(a) Laser (b) Solid state lasers (c) Semiconductor laser
Q.29 What is holography? What is the difference between holography and photography?
Q.30 What is the fundamental principle of a hologram? How is it produced and how is the image constructed
from it?
Q.31 Give briefly the requirements for holography and mention the various properties of a hologram.
Fibre Optics 5
Learning Objectives
L01 Understand the concept of optical fibre L04 Explain fibre optic communication
L02 Know about types of optical fibres L05 Illustrate optical fibre sensors,
L03 Learn about acceptance angle, connectors, and couplers
numerical aperture, skip distance and L06 Discuss the applications of optical fibre
relative refractive index couplers
Introduction
In communication systems, there has been a frequent use of either the radiowaves or the microwaves
in the form of carrier waves for sending the information. However, the advent of the laser in 1960
revolutionised the telecommunication and networking areas with an immediate appreciation of the
potential benefits of sending information from one place to the other using light, as the laser is a coherent
source of light waves. It is worth mentioning that at higher optical frequencies (~ 1015 Hz), one hundred
thousand times more information can be carried compared to microwaves. However, the energy of light
waves gets dissipated in open atmosphere. So it cannot travel long distances and hence a guiding channel
is required to guide them just like a metal wire is required to guide electrical currents. This purpose is
solved with the use of optical fibre. Optical fibre is a very thin glass or plastic conduit designed to guide
light waves along the length of the fibre. As long as the refractive index of this fibre is greater than that of
its surrounding medium, the light shall suffer a large number of total internal reflections and hence much
of the light launched into one end will emerge from the other end due to small losses.
Fibre optics is a technology that uses glass, plastic, threads or fibres to transit data. A fibre optic cable
consists of a bundle of glass threads (Fig. 5.1) which are protected by the cable’s outer covering of treated
paper, PVC or metal, called a jacket. Optical fibre has a number of advantages over the copper wire used
to make connections electrically. For example, optical fibre, being made of glass or sometimes plastic, is
protected from electromagnetic interference such as is caused by thunderstorms. A single optical fibre
has its parts as core, cladding and sheath (protecting layer), as shown in Fig. 5.2. Core is thin glass cen-
Fibre Optics 187
tre of the fibre where the light travels. Clad-

1 1 0 0 0 1 1 0 1 1
ding is outer optical material surrounding the 0 0
1 1 Fibre Optic
core that reflects the light back into the core 0 Cable
0
because cladding has lower refractive index. 0 1
Sheath is plastic coating that protects the fi-

bre from damage and moisture.
In order to understand the advantages of fi-
Metal Cable
bre optics, it is necessary to know about the
Figure 5.1
bandwidth in general. Bandwidth is the differ-
ence between the upper and lower cutoff frequencies of a filter, a com- Sheath
munication channel, or a signal spectrum. It is typically measured in Hertz. Cladding
In the case of a lowpass filter or baseband signal, the bandwidth is equal to
its upper cutoff frequency. In radio communications, bandwidth is the range
of frequencies occupied by a modulated carrier wave. For example, an FM
radio receiver’s tuner spans a limited range of frequencies. In optics, it is the
width of an individual spectral line or the entire spectral range. Core
Fibre optics has many advantages compared with traditional metal commu- Figure 5.2
nications lines, which are listed as follows:

(i) Fibre optic cables can carry more data as their bandwidth is greater than metal cables.
(ii) Fibre optic cables are less susceptible than metal cables to interference.
(iii) Fibre optic cables are much thinner and lighter than metal wires.
(iv) Through fibre optic cables the data can be transmitted digitally rather than analogically.
(v) Attenuation through fibre optic cables is very low in transmitting the data over a long distance,
so there is so need of repeaters.
5.1 FundamentaL Ideas about optIcaL FIbre LO1

Optical fibres use light to carry digital signals and as mentioned earlier the base of this technology is the
concept of total internal reflection. The digital signal that is carried by the light is reflected inside the optical
cable and hence transfers the information. Main concepts of physics that are involved in optical fibres are
refraction, refractive indices, critical angle and total internal reflection. In refraction, the light wave bends
away from the normal when it propagates from a higher refractive index medium to a lower refractive index
medium. The phenomenon of total internal reflection takes place when the angle of refraction becomes 90°.
The incident angle at which the angle of refraction (transmitting) is equal to 90° is called critical angle. When
a light wave propagating from a higher refractive index medium to a lower refractive index medium has a
sufficiently large angle, i.e., greater than the critical angle, the light gets reflected back into the same medium.
For a particular case of an optical fibre whose core is made of glass which is bounded by a plastic cladding,
the critical angle is 82°. Therefore, the light when hits the plastic cladding at an angle more than 82° would
be reflected back in the same medium, i.e., back to the glass core. This is shown in Fig. 5.3.
> 82°
Beams of
Light
Optical Fibre
Figure 5.3
5.2 optIcaL FIbres as a dIeLectrIc WaveguIde LO2

An optical fibre is a dielectric waveguide with a very high bandwidth. It guides electromagnetic waves in an
optical spectrum, the same way as microwaves are guided by rectangular or cylindrical metallic waveguides.
An optical fibre confines the propagating waves inside it by utilizing the property of total internal reflection
of light from a dielectric interface (i.e., the interface between two dielectric materials), whereas the waves
in waveguides are confined within these structures by the reflection of the waves from the walls of the
waveguides. An optical waveguide has a circular cross section, which is made in such a way that the outer
dielectric (cladding) has lower refractive index than that of the inner one (core). Due to different refractive
indices, the phenomenon of total internal reflection takes place in the optical fibre.
The transmission of light in optical fibres over large distances is possible with minimum loss of data. These fi-
bres occupy less space and are light in weight in comparison with the waveguides or transmission lines. More-
over, the optical fibre has greater tensile strength. However, optical fibres are costly and also, they need more
protection than waveguides or transmission lines. The other disadvantage is that attenuation of signals takes
place in optical fibres. The attenuation is mainly due to Rayleigh scattering (which is inversely proportional to
the fourth power of wavelength), absorption due to impurities and radiation of light due to bending of the core.
5.3 types oF optIcaL FIbres LO2

Optical fibres are categorised based on their transmission properties and the structure. These can be classified
into two types, one of which is single mode fibre and the second one is multimode fibre. The core size is the
basic structural difference in the optical fibres.
125 125
5.3.1 single mode step Index Fibre
A single mode fibre is called single (mono) mode step index fi-
bre because the refractive index of the fibre “steps” up as we
move from the cladding to the core and this fibre allows single
mode to propagate at a time due to very small diameter of its core
(Fig. 5.4a). In this fibre, the refractive indices of the cladding and
the core remain constant. The size of its core (diameter) is typi-
10 50
cally around 10 mm. Single mode fibres have a lower signal loss
Single Mode Fibre Multiple Mode Fibre
and a higher information capacity or bandwidth than multimode
fibres (introduced later) as the signal loss depends on the opera- (a) (b)
tional wavelength. These fibres are capable of transferring higher Figure 5.4
Fibre Optics 189
amount of data due to low fibre dispersion. In these fibres, the wavelength can increase or decrease the losses
caused by fibre bending. In general, single mode fibres are considered to be low loss fibres, which increase
system bandwidth and length. So these fibres are most useful for large bandwidth applications. Since these
fibres are more resistant to attenuation, they can also be used in significantly longer cable runs.
5.3.2 multimode Fibres
As the name implies multimode fibres allow more than one mode to propagate. Over 100 modes can propagate
through multimode fibres at a time. Multimode fibre is sometimes abbreviated as MMF. The size of its core is typ-
ically around 50 mm (Fig. 5.4b). The multimode fibre is of two types, namely step index and graded index fibres.
5.3.2.1 Multimode Step Index Fibres
Multimode step index fibre is shown in Fig. 5.5 along with the refractive indices of its core and cladding. In
this type of optical fibre, the number of propagating modes depends on the ratio of core diameter and the wave-
length. This ratio is inversely proportional to the numerical aperture (abbreviated as NA and defined later).
Typically the core diameter is 50 mm to 100 mm and NA varies from 0.20 to 0.29, respectively. Multimode
fibre is used in short lengths, such as those used in Local Area Networks (LANs) and Storage Area Networks
(SANs). Because the multimode optical fibre has higher NA and the large core size, fibre connections and
launching of light has become very easy. Multimode fibres permit the use of light emitting diodes (LEDs). In
such fibres, core-to-core alignment is less critical during fibre splicing. However, due to several modes the ef-
fect of dispersion gets increased, i.e., the modes arrive at the fibre end at slightly different times and so spread-
ing of pulses takes place. This dispersion of the modes affects the system bandwidth. Therefore, the core diam-
eter, NA, and index profile properties of multimode fibres are optimized to maximize the system bandwidth.
Air
m2 Cladding m2
m2
m1 m1
m1 Core
m2 Cladding m2
Air
Figure 5.5
5.3.2.2 Multimode Graded Index Fibres

In a multimode graded index optical fibre, the refractive index of the core decreases with increasing radial
distance from the fibre axis, which is the imaginary central axis running along the length of the fibre (Fig. 5.6).
The value of the refractive index is highest at the centre of the core and decreases to a value at the edge of the
core that equals the refractive index of the cladding. Therefore, the light waves in the outer zones of the core
Core m2 Cladding
m1
m
m(r) m1
m(r)
Cladding Cladding
(a) (b) Core (c)
Figure 5.6
travel faster than those in the centre of the core. Thus the dispersion of the modes is compensated by this type
of fibre design. Under this situation, the light waves follow sinusoidal paths along the fibre. In such fibres,
the most common profile of the refractive index is very nearly parabolic that results in continual refocusing
of the rays in the core, and minimizing modal dispersion. Standard graded index fibres typically have a core
diameter of 50 mm or 62.5 mm and a cladding diameter of 125 mm. It is typically used for transmitting the
information to the distances of a couple of kilometers. The advantage of the graded index fibre in comparison
with multimode step index fibre is the considerable decrease in modal dispersion.
5.4 acceptance angLe and numerIcaL aperture LO3

In order to propagate or transmit the light wave through the optical fibre, it is necessary to launch the light
at angles that fall within certain range. The maximum limit of this angle is decided by the acceptance angle.
Refracted ray Normal

B
Air m 0 q Totally internally

A qr q reflected ray
Axis of
qi C the fibre
Core m 1
tr ay
den
Inci
Cladding m 2
Figure 5.7
5.4.1 Acceptance Angle
Let us consider an optical fibre into which the light is incident. In Fig. 5.7, we show a section of cylindrical
optical fibre. The refractive index of the core is m1 and that of the cladding is m2 such that m1 > m2. The refractive
index of the medium form which the light is incident in the fibre is m0. A light wave enters the fibre at an angle
qi with the axis of the fibre. This wave gets refracted at an angle qr and strikes core-cladding interface at an
angle q. If q is greater than the critical angle qc, the wave undergoes total internal reflection at the interface,
since m1 > m2. As long as the angle q is greater than qc, the light will stay within the core of the fibre.
Let us now compute the incident angle qi for which qi ≥ qc such that the light refocuses within the core of the
fibre. Applying Snell’s law to the launching face of the fibre, we get
sin q i m1
= (i)
sin q r m0
If qi is increased beyond the limit, q will drop below the critical value qc (as qr + q = 90°, in DABC) and the
ray escapes from the side walls of the fibre. The largest value of qi occurs when q = qc. This value of qi we
represent by qi max. From the DABC, it is seen that
sin qr = sin(90° – q) = cos q, (as qr + q = 90°) (ii)
From Eqs. (i) and (ii), we get
sin q i m1 m
= or sin q i = 1 cos q
sin q r m0 m0
Fibre Optics 191
m1
when q = qc, sin qi = cos qc (iii)
m0
m2
At critical angle, sin qc = (as q = 90°)
m1
\ cos q c = 1 - sin 2 q c
m 22
= 1-
m12
m12 - m 22
or cos qc = (iv)
m12
By putting the value of cos qc from Eq. (iv) into Eq. (iii), we get
m12 - m 22
sin qi max = (v)
m0
If the incident wave of light is launched from air medium (for which m0 = 1), then
putting qi max = q0, Eq. (v) may be simplified to
sin q 0 = m12 - m 22 or q 0 = sin -1 ( m12 - m 22 ) (vi)
A Ls B
Acceptance qr
Cone q
O d
q0 qr
q i max
Figure 5.8
The angle q0 is called the acceptance angle of the fibre, which may be defined as the maximum angle that
a light wave can have relative to the axis of the fibre for its propagation through the fibre. The light wave
contained within the cone having a full angle 2q0 are accepted and transmitted along the fibre. Therefore, the
cone associated with the angle 2q0 is called the acceptance cone (Fig. 5.8). The light incident at an angle
beyond q0 refracts through the cladding. As at every internal reflection the light will be lost being incident
at an angle less than the critical angle, the corresponding optical energy is lost. It is also obvious that the
acceptance angle would be larger if the diameter of the cone is larger.
5.4.2 Numerical Aperture
Numerical aperture (NA) is the most important parameter of an optical fibre. It is a measure of how much
light can be collected by an optical system such as an optical fibre or a microscope lens. Based on the
refractive indices of core and cladding, we can measure the values of NA. It is defined as the sine of the
acceptance angle if the end faces of the fibre are exposed to a medium for which m0 = 1 (air). Otherwise, the
numerical aperture is defined as NA = m0 sin q0
For m0 = 1, NA = sin q0 = m12 - m 22

This relation shows that the light gathering ability of an optical fibre increases with its numerical aperture.
Since the maximum value of sin q0 can be 1 only, the value of NA cannot exceed 1. It means the largest value
of NA is unity. When q0 ª 90°, the fibre totally reflects all the light entering its face. Fibres with a wide variety
or numerical apertures running from about 0.2 up to 1.0 and including 1.0 may commercially be obtained.
5.4.3 Skip Distance
It is well known that the light propagates in the optical fibre based on the principle of total internal reflection.
The light ray gets reflected from the walls of the fibre. The distance between the two successive reflections
of a ray of light propagating in the fibre is called the skip distance Ls. In Fig. 5.8, the distance AB is the skip
distance, given by
Ls = d cot qr,
where d is the diameter of the core of the fibre and qr is the angle of refraction in the core. We can write
the above relation in terms of incidence angle qi and the refractive indices m1 and m0 by using Snell’s law as
sin q i m1
= . This gives
sin q r m0
m m2
sin q r = 0 sin q i or cos q r = 1 - 02 sin 2 q i
m1 m1
2
Ê m 2 sin 2 q ˆ
1- Á 0 2 i ˜
Ë m1 ¯
Hence Ls = d
m0 sin q i
m1
2
Ê m1 ˆ
or Ls = d Á -1
Ë m0 sin q i ˜¯
It is clear that the inverse of the skip distance Ls, i.e., 1/Ls will give the total number of reflections made by
the light ray in a given length of the fibre. For example, in a fibre of length L, the number of reflections Nr
would be
L
Nr =
2
Ê m1 ˆ
d Á -1
Ë m0 sin q i ˜¯
For example, in the case m1 = 1.60, m0 = 1, qi = 30° and d = 0.05 mm, we get the skip distance as 0.152 mm.
Therefore, in 1 m of fibre there will be 6580 reflections.
5.4.4 Relative Refractive Index

It has been established that the refractive indices of the core and the cladding of the fibre are different. The
difference of these two indices gives a measure of the relative refractive index difference. In this light, it can
be obtained that
Fibre Optics 193
( m1 + m 2 ) ( m1 - m 2 )
m12 - m 22 = ( m1 + m 2 )( m1 - m 2 ) = 2m1
2 m1
( m1 + m 2 ) (m - m2 )
Now is very nearly equal to m1 in view of m1 > > m2. Further, we put 1 = Dmr in the above
2 m1
equation and obtain m12 - m 22 = 2 m12 Dm r . In term of this relation the numerical aperture NA can be written as
NA = m12 - m 22 = m1 2 Dm r
Here Dmr is called the relative refractive index difference or fractional refractive index.
5.5 FIbre optIcs communIcatIon LO4

A general communication system was oftenly used before the development of optical fibre communication.
This system employs its essential components as modulator or transmitter, transmission medium and
the demodulator or receiver, as shown in Fig. 5.9a. However, optical fibres have replaced copper coaxial
cables due to their various advantages as they are very light weight, thin conduit cables which provide
greater communications capacity with lower loss.
Message Modulator Transmission Demodulator Destination

Input (Transmitter) Medium (Receiver)
(a)
Optical
Message Modulator/ Optical Fibre Optical Demodulator Destination
Input Transmitter Source Detector (Receiver)
(b)
Figure 5.9
A fibre optic communication system from signal source to signal output is shown in Fig. 5.9b. Here, the
information that is to be transmitted is first converted into an optical signal from an electrical signal. Then
the optical signal is converted to an electrical signal after transmission by an optical fibre. Independent of the
original nature of the signal, a fibre provides the choice of format of transmission as analog or digital because
these two formats are convertible into one another. So the signal in analog or digital form is impressed
onto the carrier wave by using a modulator. The carrier wave is generated from the carrier source which
may be either light emitting diode (LED) or laser diode (LD). This carrier wave is modulated using various
techniques viz., frequency modulation, amplitude modulation and digital modulation. The carrier source
output into the optical fibre is represented by a single pulse. When a pulse is passed through a fibre, then
it is attenuated and distorted due to several mechanism for example by intermodal distortion. Therefore,
repeaters and regenerators are used to amplify the light signal at several positions of the fibre. And after
that the light signal is coupled into a detector that may be a semiconductor device or most commonly
a PIN diode at the end of a fibre. This changes the optical signal back into an electrical signal. The
response of a detector should be well matched with the optical frequency of the signal received. The output
of the detector then passes through the signal processor, which is used to capture the original electrical signal
from the carrier by using the process of filtering, amplification and an analog to digital conversion. The signal
output is finally communicated by the cathode ray tube (if it is video signal), by loudspeaker (if it is audio
signal) or by computer input (if it is digital signal).
5.5.1 High Bit Rate Optical Fibre Communication

A high bit rate signal carried on copper wire transmission line is generally needed to be amplified every 300
m. However, high bit rate signals when carried on optical fibres need such amplification only every 100 km
or so. As discussed earlier, a detector changes the optical signal back into an electrical signal, the light signal
is coupled to the detector changes the optical signal back into an electrical signal, the light signal is coupled
to the detector at the remote end of the fibre. This is done effectively when the response of the detector is
well matched with the optical frequency of the signal received. Then a signal processor handles the detector
output. The function of the signal processor is to recapture the original electrical signal from the carrier. This
process involves filteration and amplification and a digital to analog conversion.
5.5.2 Allowed Modes and Normalised Frequency

It appears from the theory of acceptance cone that every ray shall propagate successfully once it enters the
fibre within its acceptance cone. However, this is not the case always and only certain ray directions or modes
are allowed to propagate successfully. Actually any ray represents plane waves that move up and down in
the fibre. Evidently such waves overlap and interfere with one another. Only those waves will sustain which
satisfy a condition of resonance. Keeping in view this point, we can derive a relation for a parameter mm in
terms of the core diameter d, numerical aperture NA and the wavelength l. This is given by
2
1Ê d ˆ
mm = Á p NA˜
2Ë l ¯
The largest integer that is less than the parameter mm shall give the maximum number of modes that propagate
successfully in the fibre. Therefore, it is clear that number of possible modes will be larger for the higher
ratio d/l. So, larger diameter fibres shall allow more number of modes to propagate. For this reason, they are
called multimode fibres. However, if d/l is small such that mm is less than 2, the fibre will allow only one
mode. So this type of fibre is called single mode fibre or monomode fibre. The condition mm < 2 for a single
mode fibre can be achieved if
d 2
<
l p ( NA)
The above condition related to the diameter guarantees the performance of single mode fibre. However, a
more careful analysis reveals that the single mode performance can be achieved even if
d 2.4
<
l p ( NA)
As is evident, the parameter mm decides the number of possible modes. Since this parameter depends on
core diameter d and the numerical aperture NA, the number of allowed modes would be different for fibres
of different core diameters. The word ‘‘number’’ intuitively adds a concept of normalised frequency, given by
d pd
nn = p NA = m12 - m 22
l l
Fibre Optics 195
A careful look indicates that the normalised frequency is nothing but the factor carried by the parenthesis of
the parameter mm. Therefore, in terms of normalised frequency nn, the parameter mm is written as
n n2
mm =
2
5.5.3 Attenuation
When light travels along the fibre, there is a loss of optical power, which is called attenuation. Signal
attenuation is defined as the ratio of optical input power (Pi) to the output power (P0). Optical input power is
the power transmitted into the fibre from an optical source. Optical output power is the power received at the
fibre end. The following relation defines the signal attenuation or absorption coefficient in terms of length L
of the fibre.
10 P
a = log10 i
L P0
So signal attenuation is a log relationship. Length L of the fibre is expressed in kilometers. In view of this,
the unit of attenuation is decibles/kilometre i.e., dB/km. The causes of attenuation in an optical fibre are
absorption, scattering and bending losses. Each mechanism of loss is influenced by the properties of fibre
material and fibre structure. However, loss is also present at fibre connections. Absorption losses over a
length L of fibre can be described by the usual exponential law for light intensity (or irradiance) I
I = I0e–aL
where I0 is the initial intensity or the irradiance of the light.
The attenuation profile for a single mode cable is 5
depicted in Fig. 5.10, which shows that the amount Attenuation / km - dB
of attenuation is also wavelength dependent. In 4
the figure, two absorption peaks at 1.0mm and 1.4
mm are observed which are respectively due to the 3
Single-mode Fibre
peculiarities of the single mode fibre and the traces
of water remaining in the fibre as an impurity. 2
The wavelengths 1.31 mm and 1.55 mm are the
two standard single mode wavelengths that are 1
commonly used due to this water absorption peaks.
However, now the wavelength 1.55 mm are used in 0.7 0.8 0.9 1.0 1.1 1.2 1.3 1.4 1.5 1.6
view of the need to extend the distance between Wavelength (mm)
repeaters.
Figure 5.10
5.4.4 Pulse Dispersion in Optical Fibre
The spreading of pulses of light as they propagate along a fibre is called dispersion. In optics, dispersion is
the phenomenon in which the phase velocity of a wave depends on its frequency. Such medium is called a
dispersive medium. The dispersive effects in a single mode fibre are much smaller than a multimode fibre.
Due to dispersion, optical pulses in optical fibres spread and hence the signals degrade over long distances.
There are several factors that cause dispersion in optical fibres. For example, in multimode fibres, different
axial speeds of different transverse modes cause intermodal dispersion that limits the performance of the fibre.
In single mode fibres, though intermodal dispersion is eliminated, chromatic dispersion occurs because of the
slight variation in the index of the glass with the wavelength of the light. Dispersion limits the bandwidth of
the fibre because the spreading optical pulses limit the rate that pulses can follow one another on the fibre and
still remain distinguishable at the receiver.
5.6 optIcaL FIbre sensors LO5

Optical fibre sensors are fibre based devices that are used for sensing some typical quantities like temperature
or mechanical strain. These sensors are also sometimes used for several vibrations, pressure, acceleration,
or concentrations of chemical species. The general principle of operation of fibre optic sensors is that when
a light beam is sent through an optical fibre, then its parameters either in the fibre or in one or several fibre
Bragg gratings experience subtle change. Then the light reaches a detector arrangement that measures these
changes (Fig. 5.11). The light beam may be changed in five of its optical properties viz. intensity, phase,
polarisation, wavelength and spectral distribution.
Optical fibre sensors have a number of advantages over other types of sensors.
(i) They consist of electrically insulating material, which makes possible their use in high voltage
environments.
(ii) Since there is no risk of electrical sparks, even in the case of defects, these can be safely used in
explosive environments.
(iii) They are immune to electromagnetic interference (EMI).
(iv) Their materials can be chemically inactive.
(v) They can operate over a broad range of temperature.
(vi) They have multiplexing capabilities, i.e., multiple sensors in a single fibre can be interrogated with
a single optical source.
Feed Return
Optical Fibre Measurement Fibre Optical
Source Zone Detector
Figure 5.11
There are two types of sensors named intrinsic sensor and extrinsic sensor, which are dicussed below.
5.6.1 Intrinsic Sensors
In these types of sensors, the sensing medium is itself a fibre. It means the propagating light never leaves
the fibre and is altered in some way by an external phenomenon. The simplest type of sensor called intensity
based fibre optic pressure sensor is based on the variation of intensity, as in this case only a simple source
and detector are required. A special feature of intrinsic fibre optic sensors is that they can provide distributed
sensing over distances of up to one metre.
This type of sensor is useful in measuring the force being exerted between the two objects A and B, shown
in Fig. 5.12. The fibre will become slightly deformed when the pressure is increased and it experiences
increased microbending losses which results in a decrease in the light intensity received at the detector. A
decrease in the pressure relieves stress on the fibre and hence there is an increase in transmitted light detected.
Fibre Optics 197
Pressure
A
Light Fibre Light
Source Detector
B
Pressure
Figure 5.12
5.6.2 Extrinsic Sensors
In extrinsic sensors, the delivery of light and its collection is done by the fibre. Thus the propagating light
leaves the fibre, is altered in some way, and is collected by the same or another fibre. These sensors are used to
measure vibration, rotation, displacement, velocity, acceleration, torque, and twisting. A major benefit of these
sensors is their ability to reach places which are otherwise inaccessible. For example, the temperature inside
aircraft jet engines is measured by using a fibre that transmits radiation into a radiation pyrometer located
outside the engine. The same way, extrinsic sensors can also be used to measure the internal temperature of
electrical transformers, where the extreme electromagnetic fields present make other measurement techniques
impossible.
An example of an intensity based extrinsic sensor is shown in Fig. 5.13, which detects any increase or decrease in
the length/between the two fibres. The amount of light launched into the return fibre will decrease as the distance
between the two fibres is increased. However, if the length is decreased the light intensity collected by the receiver
will increase. This way these fibre optic sensors are capable of determining small shifts between objects.
Light Light
Source Detector
Feed Fibre l Return Fibre

Figure 5.13
5.7 optIcaL FIbre connector LO5

The end of an optical fibre is terminated by using an optical fibre connector, which enables quicker connection
and disconnection than splicing. The coupling and alignment of the cones of fibres is done with the help of
connector so that light can pass. Now a days a variety of optical fibre connectors is available. In general,
the connectors are differentiated by their dimensions and methods of mechanical coupling. For example,
some organisations standardize one kind of connector, depending on what equipment they commonly use,
or per type of fibre. Now-a-days, small form factor connectors (for example, LC) and multifibre connectors
(for example, MTP) are replacing the traditional connectors (for example, SC) in view of their datacom and
telecom applications.
5.8 optIcaL FIbre coupLers LO5

Optical fibre couplers are fibre devices that are used for coupling light from one or several input fibres to one
or several output fibres. Optical fibre couplers can distribute the optical signal (power) from one fibre among
two or more fibres. So the light from an input fibre can appear at one or more outputs, with the distribution
of power depending on the wavelength and polarisation of

light. A two-by-two fibre coupler is shown in Fig. 5.14.
Simple point to point connections are required for some
fibre optic data links, which need multiport or other types Figure 5.14
of connections. As the input signal is divided among the
output ports, fibre optic couplers attenuate the signal much more than a connector or splice. Fibre optic
couplers can be either active or inactive devices. An inactive coupler redistributes the optical signal without
optical-to-electrical conversion, whereas the active couplers are electronic devices that split or combine the
signal electrically and use fibre optic detectors and sources for input and output.
A basic fibre optic coupler is shown in Fig. 5.15 with N1 input ports and N2 output ports, which range from 1
to 64. The number of input ports and output ports vary depending on the intended application for the coupler.
Types of fibre optic couplers include optical splitters, optical combiners, X couplers, star couplers, and tree
couplers. Fibre couplers are usually directional couplers, which means that essentially no optical power sent
into some input port can go back into one of the input ports.
N1 INPUT PORTS N2 OUTPUT PORTS
1 1
2 PASSIVE 2
3 FIBRE OPTIC 3
COUPLER
N1 N2
Figure 5.15
5.9 appLIcatIons oF optIcaL FIbre coupLers LO6

Below we mention some typical applications of fibre couplers.
(i) Fibre couplers can be used in fibre interferometers.
(ii) In a cable TV system, the powerful signal from one transmitter is sent into a fibre splitter, which
distributes the power signal over a large number of output fibres for different customers.
(iii) In fibre ring lasers, there is no resonator ends where light could be injected so a dichroic fibre
coupler can be used to inject pump light within resonator. Then another fibre coupler is used as the
output coupler.
(iv) For combining the radiation of several laser dioide in high power fibre lasers, we generally use
multimode fibre couplers.
s UmmarY
The essence of the topics covered in this chapter is produced below.

✦ In communication systems, radiowaves or microwaves have been extensively used as the carrier waves
for transmitting information. However, invention of laser led to the discovery of optical fibres which
can carry hundred thousand times more information.
Fibre Optics 199
✦ Advantages of optical fibres over the traditional metal communication lines were discussed in view
of their greater bandwidth, less susceptibility to interference, light weight, smaller thickness, and fast
transmission of data.
✦ Based on transmission properties and the structure, we can categorize optical fibres as single mode fibre
or multimode fibre. Typical diameter of core of the single mode fibre is 10 mm and that of multimode
fibre ranges from 50 mm to 100 mm.
✦ Since the refractive index steps up when we move towards core side from the cladding side, these fibres
are referred to as step index fibres.
✦ In order to compensate the mode dispersion, fibre is designed such that the refractive index of core and
cladding match at their common boundary. Such type of fibre is called multimode graded index fibre
where most commonly parabolic profile of the refractive index is used.
✦ It is not necessary that all the incident light rays shall transmit through the fibre. In this context,
acceptance angle is an important parameter. The rays that fall within the acceptance cone are accepted
for the transmission.
✦ Numerical aperture (NA) is the most important parameter of an optical fibre, which tells us how much
light can be collected by an optical fibre.
✦ It is well known that the propagation of the light is based on its total internal reflection. The light gets
reflected from the walls of the fibre. The distance between the two successive reflections of a light ray
propagating in the fibre is called the skip distance. Inverse of this distance gives the total number of
reflections made by the light in the fibre of a given length.
✦ Propagation mechanism of the information in optical fibres was discussed in detail along with a
difference of components used in the general communication system.
✦ The acceptance cone only accepts the rays for their transmission in the fibre. It is not that the every
ray shall propagate successfully once it enters within the acceptance cone. Only certain ray directions
or modes are allowed to propagate successfully. Since ray represents plane waves that move up and
down in the fibre, such waves overlap and interfere with one another. Only those waves will sustain
which satisfy a condition of resonance. Such waves or modes are called allowed modes. Therefore, the
concept of allowed modes and the normalised frequency were given and in support some theoretical
relations were talked about.
✦ When light travels along the fibre, there is a loss of optical power. This is called attenuation. Signal
attenuation is defined as the ratio of optical input power (Pi) to the optical output power (P0) and is
given by
10 P
a= log10 i
L P0
✦ The unit of attenuation is decibels/kilometre, i.e., dB/km.
✦ In addition to the loss of power of the signal (pulse of light) that propagates in the fibre, there is
spreading of pulses of light. This is called dispersion. So the dispersion was talked about in the case of
optical fibre.
✦ The wonderful application of the optical fibres is in fibre optic sensors, which are fibre based devices
that are used for sensing some typical quantities like temperature or mechanical strain. These sensors
are also sometimes used for sensing vibrations, pressure, acceleration, or concentrations of chemical
species.
✦ The types of optical fibre sensors, namely intrinsic sensors and extrinsic sensors, were discussed. In the
intrinsic sensors, the sensing medium is itself a fibre. So the propagating light never leaves the fibre and
is altered by an external phenomenon. On the other hand, in the extrinsic sensors, the delivery of light
and its collection is done by the fibre. Thus the propagating light leaves the fibre, is altered in some
way, and is collected by the same or another fibre. The extrinsic sensors are used to measure vibration,
rotation, displacement, velocity, acceleration, torque, and twisting.
✦ Another application of fibres is in optical fibre connectors and couplers. Optical fibre couplers are
fibre devices that are used for coupling light from one or several input fibres to one or several output
fibres. Optical fibre couplers can distribute the optical signal (power) from one fibre among two or
more fibres. The fibre couplers also have applications in fibre interferometers, cable TV system, fibre
ring lasers, etc.
s olVeD e XamPles
E xamplE 1 The refractive indices for core and cladding for a step index fibre are 1.52 and 1.41 respectively.
Calculate (i) critical angle (ii) numerical aperture and (iii) the maximum incidence angle.
Solution Given mcore = m1 = 1.52, mclad = m2 = 1.41
Êm ˆ
Critical angle (qc) = sin -1 Á 2 ˜
Ë m1 ¯
Numerical aperture (NA) = ( m12 - m 22 ) and

-1 2 2
Maximum incidence angle (q0) = sin [ ( m1 - m 2 )]
È 1.41 ˘
(i) q c = sin -1 Í = 68.06∞
Î1.52 ˙˚
or qc = 68.1°
(ii) NA = ( m12 - m 22 ) = (1.52) 2 - (1.41)2 = 0.5677
= 0.568
(iii) q0 = sin -1[ ( m12 - m 22 )] = sin -1 (1.52)2 - (1.41)2
= sin–1 [0.568]
q0 = 34.59° = 34.6°
E xamplE 2 Find out the numerical aperture and acceptance angle of an optical fibre, if the refractive indices
for core and cladding are 1.6 and 1.5, respectively.
Numerical aperture (NA) = ( m12 - m 22 )
Acceptance angle (q0) = sin -1[ ( m12 - m 22 )]

Fibre Optics 201
NA = (1.6) 2 - (1.5) 2 = 0.31 = 0.556

or NA = 0.556
q0 = sin–1 (0.556) = 33.78°
E xamplE 3 A light ray enters from air to a fibre. The refractive index of air is 1.0. The fibre has refractive
index of core is equal to 1.5 and that of cladding is 1.48. Find the critical angle, the fractional refractive index,
the acceptance angle and numerical aperture.
Solution Given mair = m0 = 1.0, mcore = m1 = 1.5, mclad = m2 = 1.48
-1 Ê m 2 ˆ
Critical angle (qc) = sin Á ˜
Ë m1 ¯
m1 - m 2
Fractional refractive index ( Dm r ) =
m1
-1 2 2
Acceptance angle (q0) = sin [ ( m1 - m 2 )]
Numercial aperture (NA) = ( m12 - m 22 )

-1 Ê m 2 ˆ -1 Ê 1.48 ˆ
qc = sin Á ˜ = sin Á = 80.63∞
Ë m1 ¯ Ë 1.50 ˜¯
or qc = 80.63°
m - m 2 1.50 - 1.48
Dm r = 1 = = 0.0133
m1 1.48
or Dmr = 1.33% of light.
NA = ( m12 - m 22 ) = (1.50) 2 - (1.48)2 = 0.244
q 0 = sin -1[ ( m12 - m 22 )] = sin -1[0.244]

\ q0 = 14.13°
E xamplE 4 Calculate the numerical aperture and acceptance angle of optical fibre of refractive indices for
core and cladding as 1.62 and 1.52, respectively.
Solution Given mcore = m1 = 1.62 and mclad = m2 = 1.52
Numerical aperture (NA) = ( m12 - m 22 ) and
Acceptance angle (q0) = sin -1[ ( m12 - m 22 )]
NA = (1.62) 2 - (1.52) 2 = [ 0.314]

NA = 0.56
q0 = sin–1 (NA) = sin–1 (0.56) = 34.06°
q0 = 34.1°
E xamplE 5 Calculate the refractive indices of the core and cladding material of a fibre from the following
data: NA = 0.22, Dmr = 0.012, where NA is numerical aperture,
mcore - mclad
Dm r =
mcore
mcore and mclad have usual meanings.
Solution Given: NA = 0.22, Dmr = 0.012, mcore = m1 and mclad = m2

Formula used are
m1 - m 2
Dm r = and NA = m12 - m 22
m1
m1 - m 2 m
Dm r = 0.012 = =1- 2
m1 m1
or m2 = 0.988 m1
NA = 0.22 = m12 - (0.988m1 ) 2

2
or 0.0484 = m1 [0.023856]
or m1 = 1.424 = mcore
and m2 = 0.988 m1 = 1.41
m2 = 1.41 = mclad
E xamplE 6 The refractive indices for core and cladding for a step index fibre of diameter 0.064 mm are
1.53 and 1.39, respectively. Calculate (i) numerical aperture of the fibre (ii) acceptance angle (iii) number of
reflections in 90 cm of fibre for a ray at the maximum incidence angle and for one at half this angle.
Solution Given d = 0.064 mm, mcore = m1 = 1.53 and mclad = m2 = 1.39
-1 2 2
Acceptance angle (q0) = sin [ ( m1 - m 2 )]
L
Number of reflections (Nr) =
2
È m1 ˘
Í m sin q ˙ - 1
d
Î 0 i˚
NA = (1.53) 2 - (1.39) 2 = 0.639

NA = 0.64
q0 = sin–1 (NA) = sin–1 (0.64) = 39.79°
q0 = 39.8°
Nr at qi = q0, then
90
Nr =
2
Ê 1.53 ˆ
0.0064 Á ˜ -1
Ë 1 ¥ 0.64 ¯
= 6476
q0
Nr at qi = , then
2
90 90
Nr = =
2 2
Ê 1.53 ˆ Ê 1.53 ˆ
-1 0.0064 Á -1
0.0064 Á
Ë 1 ¥ sin 19.9∞ ˜¯ Ë 0.34 ˜¯
= 3205.14
Nr = 3205
Fibre Optics 203
E xamplE 7 A graded index fibre has a core diameter of 0.05 mm and numerical aperture of 0.22 at a
wavelength of 8500 Å. What are the normalised frequency (nn) and number of modes guided in the core?
Solution Given d = 0.05 mm, NA = 0.22, l = 0.00085 mm
pd
Normalised frequency (nn) = NA and
l 2
1 Èp d ˘
Maximum number of modes guided or propagated (mm) = NA˙
2 ÍÎ l ˚
3.14 ¥ 0.05 ¥ 10-3 ¥ 0.22
nn =
0.85 ¥ 10-6
= 40.63
vn = 40.63
1
and mm = (vn ) 2
2
= 825.398
mm = 825
E xamplE 8 The refractive indices of core and cladding of a fibre are 1.465 and 1.460, respectively, and the
light of wavelength 1.25 mm is used. What should be the diameter of core for a single mode propagation? If
the core diameter is given as 50 mm, how many modes can propagate through the fibre?
Solution Given mcore = m1 = 1.465, mclad = m2 = 1.460 and l = 1.25 ¥ 10–6 m, d = ?
2.4 ¥ l
For single mode propagation, d < ,
p NA
2
1 Èp d ˘
Number of modes propagated (mm) = NA˙
2 ÍÎ l ˚
So, NA = (1.465) 2 - (1.460) 2 = 0.121

2.4 ¥ 1.25 ¥ 10-6
d< = 7.896 ¥ 10-6 m
3.14 ¥ 0.121
d < 7.9 mm
2
1 È 3.14 ¥ 50 ¥ 10-6 ¥ 0.121˘
mm = Í ˙
2Î 1.25 ¥ 10-6 ˚
1
= (15.197) 2
2
230.94
=
2
= 115.47
\ Number of modes
= 115
E xamplE 9 How many modes can propagate in a step-index fibre with a core diameter as 40 mm, if the
refractive indices of its core and cladding are 1.461 and 1.456, respectively, and the light of wavelength is
8500 Å?
Solution Given mcore = m1 = 1.461, mclad = m2 = 1.456, l = 0.85 ¥ 10–6 m and d = 4.0 ¥ 10–5 m.
2
1 Èp d ˘
Maximum mode propagated (mm) = NA˙
2 ÍÎ l ˚
and numerical aperture (NA) = ( m12 - m 22 )
So, NA = (1.461) 2 - (1.456) 2

NA = 0.121
2
1 È 3.14 ¥ 4.0 ¥ 10-5 ˘
mm = Í -6
¥ 0.121˙ = 159.83
2 Î 0.85 ¥ 10 ˚
mm = 159
E xamplE 10 Consider a slab waveguide made of Al Ga As having refractive indices for core and cladding as
3.6 and 3.55, respectively. Find how many modes can propagate in this waveguide if
(i) d = 5l and (ii) d = 50l?
2
1 Èp d ˘
Number of modes propagated (mm) = NA˙
2 ÎÍ 2 ˚
and numerical aperture (NA) = ( m12 - m 22 )
NA = (3.6) 2 - (3.55) 2 = 0.5979

(i) d = 5l, then
2
1 È 3.14 ¥ 5l ˘
mm = ¥ .5979˙ = 44.06
2 ÍÎ l ˚
mm = 44
(ii) d = 50l, then
2
1 È 3.14 ¥ 50l ˘
mm = ¥ 0.5979˙ = 4405.81
2 ÍÎ l ˚
mm = 4405
E xamplE 11 Find out the maximum core diameter of an optical fibre whose core and cladding have refractive
indices as 1.460 and 1.457, respectively, and which supports only one mode at 1.25 ¥ 10–6 m wavelength.
Solution Given: mcore = m1 = 1.460, mclad = m2 = 1.457 and l = 1.25 ¥ 10–6 m.
2.4l
Diameter of core (d) < and numerical aperture (NA) = ( m12 - m 22 )
p NA
So NA = (1.46) 2 - (1.457) 2 = 0.0935
2.4 ¥ 1.25 ¥ 10-6
\ d<
3.14 ¥ 0.0935
d < 10.22 mm
\ Maximum core diameter = 10.22 mm
Fibre Optics 205
E xamplE 12 A signal of power 5 mW exists just inside the entrance of 0.1 km long fibre. Calculate the
absorption coefficient of the fibre if the power inside the fibre be 1 mW.
Solution Given L = 0.1 km, Pi = 5 ¥ 10–6 W and P0 = 1 ¥ 10–6 W.
10 ÊPˆ
Absoprtion coefficient (a) = log Á i ˜
L Ë P0 ¯
Ê 10 ˆ Ê 5 ¥ 10-6 ˆ
or a = Á ˜ log10 Á ˜ = 69.89 dB/km
Ë 0.1¯ Ë 1.0 ¥ 10-6 ¯
or a = 70 dB/km
E xamplE 13 An optical fibre cable 3.0 km long is made up of three 1.0 km length spliced together. The
losses due to each length and splice are respectively 5 dB and 1.0 dB. What would be out put power if the
input power is 5 mW?
Solution Given a = 18/3 = 6 dB/km, Pi = 5 mW.
Ê 10 ˆ ÊPˆ
Q a = Á ˜ log10 Á i ˜
Ë L¯ Ë P0 ¯
aL 1 ÊPˆ
= ln Á i ˜
10 2.303 Ë P0 ¯
Ê P ˆ 6 ¥ 3 ¥ 2.303
ln Á i ˜ = = 4.1454
Ë P0 ¯ 10
Pi
= e 4.1454
P0
Pi 5 ¥ 10-3
or P0 = 4.1454 =
e 63.143
or P0 = 0.079 ¥ 10–3 W
or P0 = 0.080 mW
E xamplE 14 A step-index fibre has a core index of refraction of n1 = 1.425. The cut-off angle for light
entering the fibre from air is found to be 8.50°. (a) Calculate the numerical aperture of the fibre. (b) Find the
index of refraction of the cladding of this fibre (c) What would be the new numerical aperture and cut-off
angle if the fibre were submersed in water?
Solution
(a) The index of refraction for air n0 = nair = 1.0003.
The numerical aperture is found from the formula
NA = n0 sin q0max = (1.0003) sin (8.50°) = 0.1479
(b) The index of refraction of the cladding can be found from the numerical aperture using the formula
n12 – n22 = NA2
This gives n22 = n12 – NA2 = (1.425)2 – (0.1479)2 = 2.0088
n2 = 1.417
(c) The index of refraction for water n0 = nwater = 1.33. Since the numerical aperture is a property of the fibre and only
depends upon n1 and n2, it will not change when the medium outside the fibre changes.
The cut-off angle will change in case the numerical aperture is to be kept unaffected by a change in n0. It means
NA = 0.1479.
Using sin q0max = NA/n0, we get
q0max = sin–1(NA/n0) = sin–1(0.1479/1.33) = sin–1(0.1112) = 6.38°.
Q.1 Optical fibre communication uses carrier wave as

(a) radiowave (b) laser wave (c) ordinary light (d) microwaves.
Q.2 Optical fibre communication is based on the phenomenon of
(a) refraction (b) total internal reflection
(c) polarisation (d) diffraction.
Q.3 The loss in intensity of light in optical fibre is due to
(a) reflection (b) absorption (c) scattering (d) all of these.
Q.4 In single mode fibre the diameter of core is nearly equal to
(a) 10 mm (b) 100 mm (c) 50 mm (d) 125 mm.
Q.5 The inner most part of the optical fibre in known as
(a) core (b) cladding (c) sheath (d) optical fibre axis.
Q.6 The refractive indices of core (m1) and cladding (m2) of an optical fibre satisfy the relation
(a) m1 > m2 (b) m1 < m2 (c) m1 = m2 (d) none of them.
Q.7 By increasing the refractive index of core, the number of modes of propagation in an optical fibre cable
(a) remains unchange (b) increases
(c) decreases (d) none of these.
Q.8 In graded index optical fibre the refractive index of core is
(a) non-uniform (b) increase towards the axis of core
(c) same at core-cladding interface (d) all of these.
Q.9 In multimode step index fibre, the core diameter is of the order of
(a) 10 to 20 mm (b) 20 to 30 mm (c) 300 to 400 mm (d) 50 to 200 mm.
Q.10 The acceptance angle in terms of refractive index of core (m1) and cladding (m2), when the end face of
an optical fibre is exposed by the air is equal to
(a) cos -1 ( m12 - m 22 ) (b) sin -1 ( m12 - m 22 ) (c) sin -1 ( m12 - m 22 ) (d) sin -1 ( m 22 - m12 ).
Q.11 If m1 be the refractive index of core, m2 that of cladding and m0 of the medium meeting end face of fibre,
the value of numerical aperture (NA) can exceed 1 when
(a) m0 = 1 (b) m0 > 1 (c) m2 < m1 () m2 > m1.
Q.12 In intrinsic optical fibre sensor, the sensing medium is
(a) fibre (b) laser light (c) light detector (d) none of these.
s Hort -a NsWer Q UestioNs
Q.1 What is fibre optics?

Q.2 What do you mean by numerical aperture?
Q.3 What do you understand by core and cladding?
Q.4 Give name of various types of fibres.
Q.5 Why optical fibre communications are so important?
Q.6 Can more than one signal be propagated in single mode fibre?
Fibre Optics 207
P ractice P roblems
General Questions
Q.1 What is an optical fibre? Define and explain the terms
(a) Acceptance angle (b) Acceptance cone
(c) Numerical aperture (d) Relative refractive index difference
(e) Propagating modes and (f) Normalised frequency.
Q.2 What are single mode, multimode and graded index fibres? Also explain in detail the difference in
structures of single mode step index and multimode graded index fibre.
Q.3 Discuss the physical significance of numerical aperture. How does it depend on refractive indices of
core and cladding?
Q.4 Explain the allowed modes in an optical fibre. How are they related to normalized frequency?
Q.5 Discuss the propagation mechanisms of light waves in optical fibre.
Q.6 Describe schematically the basic elements of optical fibre communication system.
Q.7 Explain why does fraction of power of a signal get lost due to bending of fibre.
Q.8 Discuss the attenuation and dispersion of signals in optical fibre.
Q.9 What do you understand by optical fibre sensors. How many types of optical fibre sensors are
commonly used?
Q.10 How are optical fibre connectors and couplers needed in communication?
Q.11 What are the advantages of using optical fibre communication systems?
Q.12 Enumerate some applications of optical fibre communication system.
Q.13 Write a note on
(a) Fibre optics
(b) Application of optical fibre
(c) Numerical aperture and its physical significance
(d) Optical fibre sensors and couplers.
Electron Optics 6
Learning Objectives
L01 Understand specific charge of an L04 Explain Scanning Electron Microscope
electron and Thomson’s method (SEM), its principle, components and
L02 Learn about motion of an electron in applications
uniform and magnetic fields L05 Discuss working of Scanning Tunneling
L03 Illustrate electrostatic and Microscope (STM)
magnetostatic focusing
Introduction
The branch of physics which is concerned with beam of electrons and their deflection by means of electric
and magnetic fields is referred to as electron optics. Electron optics is also concerned with interference
of beam of electrons when they cross each other and their deflection when they pass through the spacing
in its submicroscopic structure. Electron optics is related to the wave properties of electrons which can
be treated based on quantum theory.
6.1 Specific charge of an eLectron LO1

It is very difficult to measure the elementary charge of an electron. However, it is easier to find its specific
charge which is the ratio between the charge and mass of the electron. In order to find the specific charge of
the electron, we analyse the deflection of an electron in a constant magnetic field. When a beam of electrons
enters perpendicular to the direction of magnetic field, electrons move in a circle, which is called its trajectory.
This is due to the force exerted by magnetic field. This force is called Lorentz force and it acts perpendicular
both to the magnetic field and the direction of beam propagation. If the mass, charge and velocity of electron
are m, q, v, respectively and it moves in the circle of radius r, then the Lorentz force in that simple case is
F = qvB
Electron Optics 209
This force compensates the centripetal force of the circular path. So we get
mv 2
= qvB
r
q v
or =
m Br
Therefore, we can determine q/m if we are able to measure v, B and r.
6.2 Determination of Specific charge of an eLectron:

LO1
thomSon’S methoD
The idea to measure the specific charge of electron was given by J. J. Thomson in 1897. The basic concept of
its working is related to the case when a beam of electrons is passed through the perpendicular electric and
magnetic fields. The beam of electrons will remain undeflected if the forces on the electrons of the beam due
to electric and magnetic fields are equal but opposite in direction.
A1 A B P
2
C P1 S
P2
E
O
Figure 6.1
The apparatus used by Thomson is shown in Fig. 6.1. This apparatus consist of a highly evacuated gas
container P. The electrons from the hot cathode C are accelerated and a beam of electrons is formed by a
potential difference V between the anodes A1 and A2. An electric field is applied perpendicular to the path of
electron beam by using two metal plates P1 and P2. Now a magnetic field is also applied perpendicular to the
plane of the paper (pointed out of the plane of the paper) on the beam of electrons at the same place where
electric field is acting. So after passing through the perpendicular electric and magnetic fields, the electron
beam strikes the screen S at point O. The screen is coated with a material that glows at the point of impact.
In the absence of electric and magnetic fields the beam remains undeflected and strikes the screen at point O.
However, in the presence of electric and magnetic fields the beam will deflect. Now, we adjust the strength
 
of E and B fields so that the beam of electrons meets the screen S at the same point O i.e., at undeflected
position. So in this case, the forces due to electric and magnetic fields balance each other.
Consider the electron of mass m and charge q which is moving with velocity v when it comes out through the

anode. The force on the electron due to electric field E is
Fe = qE

The force on the electron due to the magnetic field B is
FB = qvB
As discussed, these forces balance each other when we obtain the undeflected position.
qE = qvB
E
v=
B
The accelerating potential V determines the speed v of the electrons. Since, the potential energy of an electron
at the cathode appears as a gain in its kinetic energy at the anode when the electron beam is accelerated from
cathode to anode. Hence, we get
1 2
mv = qV
2
q 1 v2
=
m 2V
Substituting the value of speed v in the above relation we get
2
q 1 Ê Eˆ
= Á ˜
m 2V Ë B ¯
q E2
so =
m 2VB 2
It is clear from the above relation that the ratio q/m of charge to mass of the electron can be determined by
measuring all the quantities on the right side.
6.3 motion of an eLectron in Uniform eLectric anD

LO2
magnetic fieLDS
As we know the electron is a charged particle. So it experiences a force when an electric field is applied.
Once it gains energy, i.e it acquires some velocity, the Lorentz force can significantly change its direction of
motion. This is based on the fact that a magnetic force alone cannot work. Under the situation of combined
application of electric and magnetic fields, the electron attains different trajectories. Below we discuss some
cases and calculate the trajectories of the electron. These calculations finally contribute to the electric and
magnetic focusing of the electron.

6.3.1 Uniform magnetic field ( E = 0)
Let an electron move in a uniform magnetic field, which is constant in time also. The direction of the magnetic

field is parallel to the z-axis i.e., B = Bzˆ . As we know that the kinetic energy of the electron is not changed
by the application of the magnetic field but is changed by an electric field. Consequently, the speed of the
electron remains unchanged in the constant and uniform magnetic field. In this case, the equation of motion is
 
 dv  
F = ma = m = q (v ¥ B )
dt

dv  
or m = - e(v ¥ B ) (i)
dt
In the above equation we have used q = –e. In Cartesian coordinate system, the components of the above
equation can be written as
dvx eB
= - vy (iia)
dt m
Electron Optics 211
dv y eB
= vx (iib)
dt m
dvz
=0 (iic)
dt
From Eq. (iic), we obtain
vz = vz0 (iii)
This shows that the component of the velocity, which is parallel to the direction of the magnetic field, i.e.,
in z direction, is equal to a constant value vz0. It means that electron does not accelerate in z direction.
Differentiating Eqs. (iia) and (iib) w.r.t. t, we obtain the following relations
d 2vx eB dv y d 2v y eB dvx
2
=- and 2
=
dt m dt dt m dt
2
d vx Ê eB ˆ
2
d 2v y 2
Ê eB ˆ
= - Á ˜ vx and = - Á ˜ vy (iv)
dt 2 Ë m¯ dt 2 Ë m¯
Above equations describe a simple harmonic oscillator at frequency wc called cyclotron frequency or
gyromagnetic frequency. This frequency is defined as
eB
wc = (v)
m
Note that wc is a non-negative quantity or remains positive always.
Now, multiplying Eq. (iib) with i and then adding to Eq. (iia), we get a single equation that will govern the
motion of the electron in the x and y direction. This equation is produced below.
dvx dv y eB
+i = (- v y + ivx ),
dt dt m
d (vx + iv y ) eB
= i (vx + iv y ),
dt m
dV eB
or =i V (vi)
dt m
Here V = vx + ivy
Again Eq. (vi) can be written as
dV
= iw cV (vii)
dt
The solution of Eq. (vii) is given as:
V = v^0ei (w ct + f0 ) (viii)
Where v^0 and f0 are constants, that depend upon the initial conditions of the electron motion. Now solution
(viii) can be written as
vx + ivy = v^0{cos(wct + f0) + i sin (wct + f0)}
From the above Eq. x and y components of the velocity are obtained as
vx = v^0 cos (wct + f0) (ix)
vy = v^0 sin (wct + f0) (x)
and z-component of the velocity from Eq. (iii) is written as
vz = vz0
After squaring and adding Eqs (ix) and (x), we get, vx2 + v 2y = v^0 , which is the component of velocity
perpendicular to the magnetic field. Thus,
v 2 = vx2 + v 2y + vz2 = v^2 0 + vz20 = Constant (xi)

Eq. (xi) shows that the speed of the electron is unchanged in uniform magnetic field, where E = 0.
Now we want to analyse the trajectory of the electron, for which the coordinate x, y and z need to be obtained.
For this, we integrate Eqs (iii), (ix) and (x) under the limits x: x0 Æ x, y: y0 Æ y, z0 Æ z and t : 0 Æ t and for
getting
v^ 0
x = x0 + {sin (w c t + f0 ) - sin f0 }
wc (xii)
v^ 0
y = y0 - {cos (w c t + f0 ) - cos f0 } (xiii)
wc
and z = z0 + vz0t (xiv)
So r = (x0, y0 , z0) represents the initial position of the particle. On squaring and adding Eqs (xii) and (xiii),
we get the relation
2
Êv ˆ
( x - x0 ) + ( y - y 0 ) = Á ^0 ˜
2 2
(xv)
Ë wc ¯
v
This represents the equation of a circle with centre ( x0 , y 0 ) and radius ^0 . This radius is called Larmour
wc
radius rL. The coordinates x0 and y 0 are defined as
v^ 0 v^ 0
x0 = x0 - sin f0 and y 0 = y0 + cos f0
wc wc
(xvi)
From Eq. (xv), we observe that the motion of the electron lies in xy plane, which is perpendicular to the
v
magnetic field direction. The trajectory of the electron is a circle with center ( x0 , y 0 ) and radius ^0 .
wc
The direction of gyration of the electron is such that the magnetic field generated by the electron is always
opposite to the externally applied field. It can be seen from Eqs (xii) and (xiii) that the electron is moving in
the counterclockwise direction along the circumference of the circle with a uniform angular velocity wc, as
shown in Fig. 6.2.
From Eq. (xiv), we observe that the z coordinate of the electron, which is parallel to the direction of the
magnetic field, is increasing uniformly with time. Hence, the trajectory of the electron is a helix with axis
parallel to the magnetic field direction (z-axis) and passing through ( x0 , y 0 ) . The radius of the helix is equal
to the Larmour radius and the pitch is equal to 2pnz0/wc; pitch is the distance traveled by the electron in
Electron Optics 213
completing one revolution. Hence, it is clear that the parameters of the helix depend on the initial velocity
and position of the electron. If we consider another case when the perpendicular component of the velocity
is zero i.e., v^0 = 0, the trajectory of the electron motion is a straight line along the magnetic field direction.
Also if vz0 = 0, i.e., z component of velocity which is parallel to the direction of the magnetic field is zero,
the trajectory of the electron motion is a circle with centre ( x0 , y 0 ) that is called guiding centre. Generally,
guiding centre moves in the direction of the magnetic field with a constant velocity vz0 when the electron
motion is helical. Thus, we can say that the helical motion may be separated in two ways. The first one is
the uniform motion along the magnetic field and the other one is the circular motion around the field lines.
Y
ν^ 0
wc e–
~
y0
fo
y0
→
B = z^ B
~ X
x0 x0
Figure 6.2
6.3.2 Constant and Uniform Electric and Magnetic Fields

In the presence of the electric field, the motion of an electron will be the sum of two motions i.e., the circular
Larmour gyration due to the magnetic field in addition with a drift of the guiding centre due to the electric
field. For convenience, let us separate the electron velocity and electric field in two parts, which are as follows
   
E = zEˆ z = E^ and v = zv ˆ z + v^ (i)
 
Here E^ and v^ are the components which are perpendicular to the magnetic field. These are called
perpendicular components. In the presence of both electric and magnetic fields the equation of motion for a
moving electron becomes
   
dv
m = - e( E + v ¥ B ) (ii)
dt
We separate out the parallel and perpendicular components of Eq. (ii).
dvz
m = - eE z (iii)
dt
  
dv^ 
and m = - e( E^ + v^ ¥ B ) (iv)
dt
From the integration of Eq. (iii), we get

e
vz = - E z t + vz 0 (v)
m
Again integrating Eq. (v), we get the z-coordinate,
e
z=- E z t 2 + v z 0 t + z0 (vi)
2m
Here vz0 and z0 are the initial values of vz and z, respectively.
Let us consider a particular solution vE of Eq. (iv), which is a constant value. Then, by substituting vE to be
constant in (iv), we get
  
0 = E^ + vE ¥ B (vii)

Note that vE is the cross-field velocity called drift velocity. It can be assumed perpendicular to the magnetic

field without loss of generality. Now, taking the cross product of Eq. (vii) with B , we find that
    
0 = E^ ¥ B + (vE ¥ B ) ¥ B
       
0 = E^ ¥ B + B (vE ◊ B ) - vE ( B ◊ B )
  
0 = E^ ¥ B - vE B 2
 
 Ê E ¥ Bˆ
vE = Á ^ (viii)
Ë B 2 ˜¯
     
where vE ◊ B = 0 , as vE is perpendicular to B . We can see from Eq. (viii) that the E ¥ B drift is independent

of the charge (–e), mass (m) and velocity (v) of the electron. The above result verifies that vE is a particular

dv
solution of the Eq. (iv). The neglect of term ^ is also justified, as this term gives only the circular motion
dt
at wc. So, the general solution of Eq. (iv) is looked in the form
  
v^ = vE + vl (ix)

Here vl is due to the applied magnetic field. Substituting this relation, Eq. (iv) takes the form or
    
dv  
m l = - e( E^ + vE ¥ B + vl ¥ B)  vE is constant.
dt
 
dv 
or m l = - e(vl ¥ B) (x)
dt

Here vl corresponds to the orbital motion of the electron due to the application of magnetic field alone. It
is called gyration velocity. Thus, we can say that in the presence of both electric and magnetic fields, the
velocity perpendicular to the magnetic field consists of a constant drift velocity. It is also perpendicular to
the electric field direction superimposed on an orbital motion that is same as in the constant and uniform
magnetic field alone.
In order to discuss the trajectory of the electron, it is assumed that Ez = 0, vz0 = 0 and z0 = 0. So in this case,
electron motion is in the XY plane. Further, it is taken that the electric field is directed in the y-direction, i.e.,

E = yE
ˆ y . Since drift velocity (a constant velocity) satisfies the relation,
  
0 = E^ + vE ¥ B
Electron Optics 215
the sum of this drift velocity with the velocity which is obtained when only magnetic field is applied, i.e., the
gyrating motion, will satisfy the equation of motion given by Eq. (ii) .

By using the x and y components of velocity from the previous case (when E = 0) and Eq. (viii) and (ix), we
get the relations
Ey
vx = + v^0 cos (w ct + f0 ) (xi)
B
v y = v^ 0 sin (w c t + f0 ) (xii)
We can obtain the initial values of v^0 and f0 in terms of initial velocities vx0 and vy0.
1/2
È 2 ˘ vy0
v^ 0 = ÍÊÁ v - y ˆ˜ + v y 02 ˙
E
and tan f0 = (xiii)
ÎË
x 0
B¯ ˚ vx 0 - E y /B
For a particular case in which electron is initially at rest (vx0 = vy0 = 0) Eq. (xiii) results
Ey
v^ 0 = - and f0 = 0 (xiv)
B
We can obtain x and y coordinates by integrating Eqs. (xi) and (xii) and then substituting the value of v^0 and
f0. The calculated coordinates are obtained as follows
Ey v
x= t + ^0 [sin (w c t + f0 ) - sin f0 ] (xv)
B wc
v
y = - ^0 [cos (w ct + f0 ) - cos f0 ] (xvi)
wc
The trajectory defined by the above two equations is a cycloid and the coordinates of the origin are considered
to coincide with the initial position of the particle. It is possible to get different trajectories by applying
different initial conditions. Some of the trajectories with different initial conditions are discussed below.
6.3.2.1 When vx0 = vy0 = 0
 
In this case, initially the electron is at rest. Therefore, the force v ¥ B corresponding to magnetic field does
not act on the electron. However, at this time the electric field that is perpendicular to the magnetic field is
acting. So, the electric force is directed towards the negative y-direction because of the electric field direction
along the positive y-axis. So, the electron gets accelerated in the negative y-direction with the action of the
 
electric force. As the electron acquires some velocity, the v ¥ B force will start acting upon it. This will
modify the trajectory of the electron and shall force it to move in the positive x-direction. So the electron
moves in the positive X and negative y-direction, i.e., in XY plane.
From Eqs (xi) and (xiv), we obtain the relation
vx B
= 1 - cos w c t
Ey
vx B
or vx = 1 - cos w c t where = vx
Ey
Similarly, Eq. (xii) can be written as

vy B
v y = - sin w c t , where = v y
Ey
Since Ey/B has the dimensions of velocity, it follows that vx and v y are dimensionless quantities. From
Eqs. (xv) and (xiv), we obtain that x coordinate as follows
Ey Ey
x= t- sin w ct
B Bw c
xBw c
or = w ct - sin w ct
Ey
xBw c
or x = w c t - sin w ct where = x
Ey
Similarly Eq. (xvi) can be written as
yBw c
y = cos w c t - 1 where = y
Ey
Since Ey/wcB has the dimensions of length, it follows that x and y are dimensionless quantities. Finally, we
summarise these results as
vx = 1 - cos w c t v y = - sin w c t (xvii)
x = w ct - sin w ct y = cos w c t - 1 (xviii)
The above relations show that the motion of the electron is a cycloid motion between t = 0 and t = 2p/
wc. Initially, the Lorentz force does not act on the electron as it starts from rest. The electric field forces
the electron in the negative y direction. As the velocity of the electron increases the Lorentz force tends to
curve the trajectory in the positive x direction. It means the electron moves in the positive X and negative
y directions. So, vy becomes more and more negative, reaches a minimum, and finally becomes zero at t =
2p/wc after getting less and less negative. During the time 0 £ t £ p/wc, the velocity component vx is increased
to reach a maximum value. For t ≥ p/wc, the Lorentz force continues to curve the electron in anticlockwise
direction. Since vy is positive, the electron starts moving in the positive y direction. Finally vy becomes more
and more positive and reaches a maximum. Then it starts decreasing and finally zero at t = 2p/wc. During
p/wc £ t £ 2p/wc, vx gets decreased and finally reaches zero at t = 2p/wc. So at t = 2p/wc, the electron is again at
rest. This cycloid repeats itself every 2p/wc seconds. The trajectory given by Eq. (xviii) is shown in Fig. 6.3a.
6.3.2.2 When vx0 = Ey/B,vy0 = 0
In this case, we consider that the electron is moving with some initial velocity that is equivalent to drift
velocity (Ey/B). The perpendicular component of velocity v^0 is zero in this case (from Eq. xiii).
With initial conditions, Eqs (xi), (xii), (xv) and (xvi) are defined as follows:
vx = 1 v y = 0 (xix)
x = w ct y = 0 (xx)
The above relations state that the electron moves uniformly along the x-direction with the initial velocity,
which is equal to the drift velocity. In the present case, the force due to the electric field is exactly cancelled
Electron Optics 217
by the force due to the magnetic force. So, here we can say that the total force acting on the electron is zero.
Consequently, the electron initial velocity is maintained continuously. The electron trajectory from this case
is shown in Fig. 6.3b.
~
y
8 16 24 ~
D=0 0 x
–2 (a)
~ E - Field
y B - Field
8 16 24 ~
0 x
(b)
~
y
8 16 24 ~
0 x
D = 0.5
(c)
–2
~
y
–2
D = 2.5
~
x
0 8 16 24
(d)
Figure 6.3
6.3.2.3 When 0 < vx0 < Ey/B; vy0 = 0

Let us consider that the electron is moving with some initial velocity whose magnitude lies between zero and
drift velocity and is in the x-direction. So, vx0 – Ey/B < 0 as the magnitude of electric force eEy is greater than
the magnitude of magnetic/Lorentz force evx0B, and x component of velocity vx attains the minimum value
vx0.
From Eq. (xiii) we observe that
v^0 = vx0 – Ey/B
Using this, Eq. (xi) becomes
Ê
Ey E ˆ
vx = + Á vx 0 - y ˜ cos w ct
BË B¯
vx B Ê vx 0 B ˆ
=1+ Á - 1 cos w ct
or,
Ey Ë Ey ˜¯
or, vx = 1 + ( D - 1) cos w ct

vx 0 B v B
where D= and vx = x
Ey Ey
Similarly, v y = ( D - 1)sin w ct
In view of obtaining the trajectory of the electron, we find the expression for X coordinate from Eq. (xv), as
Ey 1 Ê Ey ˆ
x= t+Á vx 0 - ˜ sin w c t
B wc Ë B¯
xBw c Ê Bvx 0 ˆ
= w ct + Á - 1 sin w c t
or, Ey Ë Ey ˜¯
or, x = w ct + ( D - 1) sin w ct
xBw c vx 0 B
where x = and D=
Ey Ey
Similarly, y = (D – 1)(1 – cos wct)
So, the position coordinates and the velocity components can be reproduced as
vx = 1 + ( D - 1) cos w ct v y = ( D - 1) sin w ct (xxi)
x = w ct + ( D - 1)sin w ct y = ( D - 1)(1 - cos w ct ) (xxii)
 
Since the electron is moving with some initial velocity, magnetic force – e(v ¥ B) does not completely
cancel the electric force initially. This results the electron to move in the negative y-direction. The magnetic
force turns the orbit in the anticlockwise direction. So, the trajectory of the electron is a cycloid, the same
as in case 6.3.2.1 with a difference that minimum value of the X-component of the velocity vx is a non-zero
positive quantity. So, the electron never comes at rest but always remains moving in the positive x-direction.
The trajectory corresponding to Eq. (xxii) is shown in Fig. 6.3c.
6.3.2.4 When vx0 > Ey/B; vy0 = 0
In this case, the electron is moving initially with some velocity, which is greater than the drift velocity. The
velocities and coordinates are same as in the previous case. Since the initial velocity is greater than the drift
 
velocity in the positive x-direction, the v ¥ B force dominates over the electric force ( evx0B > eEy) that is
directed toward negative direction. Hence, magnetic force curves the electron in the anticlockwise direction
and the electron motion starts in the positive y-direction. The direction of the electric force is opposite to
the direction of electron motion. For this reason, the electron is decelerated in that direction. Consequently,
speed of the electron is larger at the bottom portion of the orbits than at the top. So, the Larmour radius is
longer at the bottom and smaller at the top, which means the radius of curvature of the trajectory is smaller
at the top and larger at the bottom. This difference in Larmour radii at the top and the bottom portions of the
trajectory results in a drift in the positive x-direction. So, the guiding center moves perpendicular to electric
and magnetic fields. The trajectory for this case is shown in Fig. 6.3d.
Electron Optics 219
6.4 eLectroStatic anD magnetoStatic focUSing LO3

As discussed in the previous section, the application of electrostatic field and magnetostatic field can control
the electron motion. This observation has led to the concept of focusing of electron beam either by using the
electrostatic field or the magnetostatic field.
ELECTROSTATIC LENS ELECTROSTATIC LENS
2nd
1st FOCAL
FOCAL POINT
POINT
GRID
SCREEN
FOCUSING ANODE ACCELERATING

ANODE
– +
FOCUS CONTROL
Figure 6.4
6.4.1 Electrostatic Focusing
The example of electrostatic focusing is an electron gun. In the electron gun, the electrons are focused by an
electrostatic field. Electrostatic lenses are formed when negative and positive fields are near to each other.
The electron gun is formed with several parts. For example, a heater and a cathode are used to generate
electrons, a control grid is used to control electron flow, and also two anodes are used. The main purpose of
the first anode is to focus the electrons into a narrow beam on the screen. Hence, it is called focusing anode.
The second anode accelerates the electrons as they pass through it. So this is called accelerating anode. The
control grid is a hollow metal tube placed over the cathode having a small opening in the center of a plate at
the end opposite to the cathode. It can control the number of electrons that are emitted because it is near the
cathode. This is based on the fact that the negative voltage of the grid can be varied either to control electron
flow or to stop it completely. The anodes consist of two cylinders that contain plates with small holes in their
centers. The cathode is indirectly heated, so it emits a cloud of electrons. The control grid is maintained at a
negative potential with respect to the cathode to keep the electrons bunched together. A high positive potential
on the anodes pulls electrons through the hole in the grid. Now, two electrostatic fields that exist between the
control grid and first anode and between the first and second anodes focus the electron beam. The motion of
the electrons through the electron gun (by dashed lines) and the relative voltage relationships on the electron
gun elements is shown in Fig. 6.4.
The cathode (K) is at a fixed positive voltage with respect to the ground. The grid is at a variable negative
voltage with respect to the cathode. A fixed positive voltage of several thousand volts is connected to the
accelerating anode. The potential of the focusing anode is less positive than the potential of the accelerating
anode. The electrostatic field areas are often referred to as lenses because the fields bend electron streams in
the same manner as the optical lenses do with light rays. The first electrostatic lens causes the electrons to
cross at the first focal point within the field. The second lens bends the spreading streams and returns them to
a new second focal point at the screen.
6.4.2 Magnetostatic Focusing
We can take an example of electron microscope in order to explain the magnetostatic focusing. Electron
microscopes have magnetic lenses that are similar to simple solenoids. A coil of copper wire produces a
magnetic field that is shaped by the surrounding iron fixture into an optimum geometry to produce the
lensing action. As an electron moves through the magnetic field, it experiences a radial inward force, which
  
is proportional to the Lorenz force, v ¥ B , where v is

the electron velocity and B is the magnetic flux density. q
The lensing action is similar to that of an optical lens, Electron Beam
in which a ray parallel to the axis of the lens is bent to B
VL
the lens axis at the focal length of the lens. In an optical
lens, the focal length is fixed by the curvature of the lens
surfaces and cannot be changed. In the electromagnetic
lens, the focal length depends on two factors: the gun

voltage, which determines the electron velocity v , and
the amount of current through the coil, which determines BR

the magnetic field, B . Therefore, the operator controls
the focal lengths of the lenses by adjusting the currents
BL B
supplied to them. An increase in current increases the
radial force experienced by the beam and thus reduces
the focal length. A typical magnetic lens is shown in Fig. Electron spiraling down the axis
6.5. The focal length f of such a lens is given by Figure 6.5
C1V
f = ,
N 2I 2
where V is the accelerating voltage, N is the number of turns in the coil, I is the current in the coil and C1 is
a constant.
For achieving good lens characteristics, it is essential to have constant accelerating voltage and constant lens
current. Like the optical lenses, these lenses also suffer from defects, namely spherical aberration, chromatic
aberration and astigmatism. Spherical aberration is caused by lens field acting inhomogeneously on the off
axis rays, i.e., inability of a lens to focus all of a parallel incident beam to a point. This is reduced by using
stop down lens. Chromatic aberration is caused by variation in the electron energy and thus the electrons are
not monochromatic. So the electrons with different energies have different wavelengths and focus at different
points. This demands for constant accelerating voltage V. Also source of electrons needs to be coherent, i.e., it
should be of narrow range of energies. Astigmatism is caused by asymmetry in lens geometry. So, additional
coils that introduce astigmatism in a controlled way are used to correct this effect.
6.5 SCanning ElECtron MiCroSCopE (SEM) LO4

Scanning Electron Microscope (SEM) is a microscope that uses electrons rather than light to form an image.
These microscopes were developed due to the limitations of Light Microscopes. The first SEM was made
available commercially around 1965. Its late development was due to the electronics involved in “scanning”
Electron Optics 221
the beam of electrons across the sample. SEM is one of the most heavily used instruments in academic/lab
research areas and industry due to the combination of higher magnification, larger depth of field, greater
resolution, and compositional and crystallographic information.
6.5.1 Sem principle
The SEM uses a focused beam of high energy electrons to generate a variety of signals at the surface of
solid specimens (Fig. 6.6). With the help of signals emitted as a result of electron specimen interaction we
gather information about the sample including external morphology (texture), chemical composition, and
crystalline structure and orientation of materials making up the sample.
Electron Gun
Electron
Beam
Anode
Magnetic
Lens To TV
Scanner
Scanning
Coils
Backscattered
Electron Secondary
Detector Electron
Detector
Stage Specimen
Figure 6.6
In most applications, data are collected over a selected area of the surface of the sample, and a 2D image
is generated that displays spatial variations in these properties. The SEM is also capable of performing
analyses of selected point locations on the sample. Depending on the incident energy of the electron beam, a
variety of electrons (auger, secondary and back scattered), X-rays (characteristic and Bremsstrahlung), light
(cathodoluminescence) and heat (phonons) are emitted (Fig. 6.7). Several of these interactions are used for
imaging, semi-quantitative analysis and/or quantitation analysis.
6.5.2 SEM Components
All the SEMs consist of a column, a specimen chamber, detectors and viewing system. A column is used to
generate a beam of electrons. The electron beam interacts with the sample in a specimen chamber. Detectors
are used to monitor the different signals that result from the electron beam and sample interaction. A viewing
system is used to build an image from the detector signal. Essential components of all SEMs are electron
source (gun), electron lenses, sample stage, detectors for all signals of interest, display/data output devices
and infrastructure requirements like power supply, vacuum system, cooling system, etc.
Incident Electron Beam
Characteristic X-rays
Backscattered Electrons
Secondary Electrons Bremsstrahlung X-rays
Visible Light
(cathodoluminescence)
Auger Electrons
Heat
Sample Surface
Diffracted Electrons
Transmitted Electrons
Figure 6.7
6.5.3 image formation in Sem

The primary electron beam is scanned over specimen in raster pattern in synchronisation with beam in
Cathode Ray Tube (CRT), as shown in Fig. 6.8. Intensity at spot (say P1) on CRT is proportional to the signal
detected from the spot P1 on the specimen which is modulated by amplifier.
Magnification The magnification is simply the ratio of the length of the scan on the Cathode Ray Tube
(CRT) to the length of the scan l on the specimen. For a CRT screen that is 20 cm square, the magnification
is given by
C 20 cm
M= =
l l
So this is clear that the increased magnification M can be achieved by decreasing l.
Image Contrast If SA and SB represent signal generated from two points of scanned area (as shown in Fig.
6.8), then the image contrast is defined as
S A - S B DS
C= =
SA SA
In order to detect small objects of poor contrast, it is necessary to use a high beam and a slow scan speed. This
improves the signal to noise ratio, i.e., S/N ratio. Secondary electrons give topographical contrast and back
scattered electrons give atomic number contrast (Fig. 6.8).
Resolution The maximum resolution that can be achieved using a microscope means the smallest
interval distinguishable between two adjacent points. Any magnification exceeding such maximum would
not make sense since further information cannot be provided. The maximum resolution mainly depends
on the wavelength of the radiation selected for the image. Beams entering the lens and aperture system of
the microscope produces overlapping diffraction patterns for each object point. The distance rm between
Electron Optics 223
Good Contrast Poor Contrast

255.00 255.00
SA P1
Signal
Noise
SB P2
0.00 0.00
0 Pixels 251 0 Pixels 251
Figure 6.8
two diffraction maxima must exceed full width half maximum (FWHM), otherwise the diffraction maxima
cannot be seen clearly (separate). This is clear from Fig. 6.9, which is prepared based on Rayleigh criterion.
According to this, the distinction is possible when the maximum of the zero order coincides with the first
minimum of the second diffraction pattern. The distance between the two first minima, i.e., r1 is inversely
proportional to the diameter of the aperture. Diffraction patterns are dependent on the wavelength l, the index
of refraction of the surrounding medium m, and the angle q formed by the optical axis and the edge beam,
which can only just pass through the aperture. In view of this rm is given as
r1 0.61l
rm = =
2 m sin q
Intensity rm
Distance
r1
r1
Figure 6.9
As discussed in previous chapter, the product m sin q is referred to as numerical aperture. In SEM, the
resolution is limited to few nm, which is due to the electron probe size that in turn depends on the quality
of the objective lens and the electron gun. Ultimate resolution obtainable in an SEM image is limited by the
minimum probe size that can generate an adequate signal at the sample.
6.5.4 applications of SEM
The SEM has got applications in various fields, some of which are mentioned below.
(i) The SEM is routinely used to generate high resolution images of shapes of objects and to show
spatial variations in chemical compositions.
(ii) This instrument is also widely used to identify phases based on qualitative chemical analysis and/or
crystalline structure.
(iii) This instrument is used in precise measurement of very small features and objects down to 50 nm in
size.
(iv) Backscattered electron images can be used for rapid
discrimination of phases in multiphase samples.
(v) SEMs equipped with diffracted backscattered electron
detectors can be used to examine microfabric and
crystallographic orientation in many materials.
(vi) Figures 6.10 and 6.11 show SEM photomicrographs of
MESFET ohmic contacts and gate contact with poor
surface morphology, respectively. Poor ohmic metal
surface morphology leads to poor contact resistance of
source drain pads (Fig. 6.10), whereas underdeveloped
gate pattern leads to poor metal semiconductor interface Figure 6.10
(Fig. 6.11).
Figure 6.11
Electron Optics 225
6.6 SCanning tUnnEling MiCroSCopE (StM) LO4

Gerd Binnig and Heinrich Rohrer invented Scanning Tunneling Microscope (STM) in early 1980s. It is a
microscopical technique that allows the investigation of electrically conducting surfaces down to atomic
scale. The basic principle of STM is the tunneling of electrons between the sharp metallic tip of a probe and
the surface of a sample under study. By etching or tearing a thin metal wire we can obtain these tips with
only one atom at the tip. The tip is mechanically connected to the scanner, an XYZ positioning device realised
by means of piezoelectric materials (Fig. 6.12). A small current called tunneling current flows if the tip is
in contact with the sample which is positively and negatively biased. This small tunneling current is then
amplified and measured. The flow of tunneling current is very sensitive to the distance between the tip and
the sample. With the help of the tunneling current the feedback electronic keeps the distance between tip
and sample constant. If the tunneling current exceeds its preset value, the distance between tip and sample is
increased and if it falls below this value, the feedback decreases the distance. For keeping the flow of current
constant, the height of the tip is continually adjusted when tip is swept over the surface. By recording the
height fluctuations of the tip accurately, a map of the “bumps” on the surface is obtained. The tip is scanned
line by line above the sample surface following the topography of the sample.
Positioning
Device for
X, Y and Z
X, Y, Z
Tunneling
Z Current
Tip Y
–
+ Computer and
Bias Feedback Electronic
Sample –
Figure 6.12
The physical behaviour of the tunneling current provides the extreme magnification capabilities of the STM
down to the atomic scale. Quantum mechanically we can explain the case of flow of tunneling current across
the small gap that separates the tip from the sample. The tunneling current I is proportional to the tunneling
bias U but it decays exponentially with an increase of the gap (d), as per the following relation
I = K1Ue - k2d ,
where K1 and k2 are constants. The variation of I with the gap d is shown in Fig. 6.13. It is clear from the
figure that a very small change in the tip sample separation induces a large change in tunneling current
(please see dd and corresponding dI). So, the tip separation is controlled very exactly and tunneling current
is carried by the outermost atom of the tip. A feedback loop constantly monitors the tunneling current and
makes adjustments to the tip to maintain a constant tunneling current. These adjustments are recorded by the
I
Tip
dI
–
Sample I = K1Ue–k d
2
d +
U d
–
dd
Figure 6.13
computer and presented as an image in the STM software. Such a setup is called a “constant current” image.
In addition, for very flat surfaces, the feedback loop can be turned off and only the current is displayed. This
is called a “constant height” image.
s UmmarY
The main topics covered in this chapter are summarised below.
✦ Electron optics is the branch of physics which is concerned with beam of electrons and their deflection
by means of electric and magnetic fields.
✦ In order to formulate the electron optics, it is necessary to analyse the motion of the electrons in electric
and magnetic fields. Therefore, we derived the expressions for the electron trajectories for the different
 
combinations of E and B fields.
✦ The concept of gyratory motion, guiding centre and cycloid motion were introduced and explained.
✦ Having made a basic background of the electron motion in the electric and magnetic fields, the focusing
of electrons was discussed.
✦ Electrostatic focusing of electron is done using static electric field. An example of the electrostatic
focusing is the electron gun. The electrostatic field areas are referred to as the lenses because the fields
bend the electron streams in the same manner as the optical lenses do.
✦ Magnetic focusing of electron is achieved using magnetic field. An example of the magnetostatic
focusing is an electron microscope. In this case, the magnetic field is designed such that the electron
experiences a force radially inward. So the lensing action is similar to that of an optical lens in which a
ray parallel to the axis of the lens is bent to the lens axis at the focal length. In the case of magnetostatic
focusing, the focal length depends on the gun voltage V and amount of current I applied to the coil
CV
having N number of turns. The focal length f of a typical magnetic lens is given by f = 21 2 , where
C is a constant. N I
1
✦ New topics on Scanning Electron Microscope (SEM) and Scanning Tunneling Microscope (STM)
were discussed in detail. SEM uses electrons rather than the light to form an image. This microscope
was developed due to the limitations of light microscopes.
Electron Optics 227
s olVeD e XamPles
E xamplE 1 The voltage across the electrodes of a cathode ray gun is 500 eV. Calculate
(i) the energy gained by electron.
(ii) the speed of the electron.
(iii) the momentum of electron
Given mass of electron = 9 ¥ 10–31 kg.
Solution Given V = 500 eV, m = 9 ¥ 10–31 kg and e = 1.6 ¥ 10–19 C.
So,
(i) The energy gained by electron = eV
= 1.6 ¥ 10–19 ¥ 500 = 8 ¥ 10–17 J
(ii) When an electron is accelerated under a potential difference V then electron acquires a speed v. So
1
K.E gained by electron = mv2 = eV
2
2 eV 2 ¥ 8 ¥ 10-17
or v= = = 1.33 ¥ 107 m/sec
m 9 ¥ 10-31
(iii) Momentum of the electron = m
= 9 × 10–31 × 1.33 × 107
= 12 × 10–24 kg/sec
E xamplE 2 What is the momentum of acceleration of an electron of speed 2.5 ¥ 106 m/sec in a magnetic field
of 2.0 G? Given that e/m = 1.76 ¥ 1011 C/kg?
Solution Given v = 2.5 ¥ 106 m/sec, B = 2.0 G = 2.0 ¥ 10–4 T
and e/m = 1.76 ¥ 1011 C/kg.
Force on electron due to magnetic field provides the required centripetal force. Therefore,
mv 2
= Bev
r
v2
Centripetal acceleration =
r
v 2 Bev
=
r m
= 2.0 ¥ 10–4 ¥ (1.76 ¥ 1011) ¥ 2.5 ¥ 106
= 8.80 ¥ 1013 m/sec2
E xamplE 3 In a Thomson’s set up for determining e/m, the same high tension d.c. supply provides potential
to the anode of accelerating column, as also to the positive deflecting plate in the region of crossed fields. If
the supply voltage is doubled, by what factor should the magnetic field is increased to keep the electron beam
undeflected?
Solution e/m of electrons in a Thomson’s experiment is given by

e E2 V
= 2
where E =
m 2VB d
e V2 V
\ = =
m 2V B 2 d 2 2 B 2 d 2
This is also the condition for undeflected beam.
Since d is fixed from the above expression we get
B2 µ V or Bµ V
So, when V is doubled (V ¢ = 2V) then B¢ will be given by
B ¢ µ 2V = 2B
It means that if V is doubled, B should be increased to 2B to keep the electron beam undeflected.
E xamplE 4 (a) A monoenergetic electron beam with electron speed of 5.20 ¥ 106 m/sec is subject to a
magnetic field of 1.30 ¥ 10–4 T, normal to the beam velocity. What is the radius of the circle traced by the
beam?
Given e/m for an electron as 1.76 ¥ 1011 C/kg.
(b) Is the formula you employed in (a) valid for calculating radius of the 210 MeV electron beam? If not, in
what way is it modified?
Solution Given v = 5.20 ¥ 106 m/sec,
B = 1.3 ¥ 10–4 T
e/m = 1.76 ¥ 1011 C/kg
Force exerted by the magnetic field on the electron
F = e | v ¥ B | = evB sin q = evB ( q = 90°)
Since, the normal magnetic field provides the centripetal force
mv 2
evB =
r
mv v
r= =
qB (e / m) B
5.20 ¥ 106
r=
1.76 ¥ 1011 ¥ 1.30 ¥ 10-4
= 0.227 m = 22.7 cm
(b) Energy of the electron beam = 20 MeV

= 20 ¥ 1.6 ¥ 10–19 ¥ 106 J
1
= 20 ¥ 10-13 ¥ 1.6 = mv 2
2
1
Ê 2 ¥ 20 ¥ 10-13 ¥ 1.6 ˆ 2
\ v=Á ˜
Ë 9 ¥ 10-31 ¯
v = 2.67 ¥ 109 m/sec
which is greater than the velocity of light.
Electron Optics 229
mv
So, the formula r = is not valid for calculating the radius of the path of 20 MeV electron beam because electron with
qB
such a high energy has velocity in the relativistic domain (comparable with the velocity of light). Since the mass varies at
such speed, we use relativistic formula as follows
mv mo Ê v ˆ
r= = Á ˜
qB 1 - v 2 / c 2 Ë qB ¯
E xamplE 5 In a Thomson set up for determination of e/m, electrons accelerated by 2.5 kV enter the region of
crossed electric and magnetic fields of strengths 3.6 ¥ 104 V/m and 1.2 ¥ 10–3 T respectively and go through
undeflected. Determine the e/m of an electron.
Solution Given V = 2.5 kV, E = 3.6 ¥ 104 V/m, and B = 1.2 ¥ 10–3 T.
Energy gained by the electron = eV
1
eV = mv 2
2
e v2
or =
m 2V
Since electrons go through undeflected from the region of crossed electric and magnetic field, we have
E
eE = Bev fi v =
B
e ( E /B ) 2 [3.6 ¥ 104 /1.2 ¥ 10-3 ]2 9 ¥ 1014
= = =
m 2V 2 ¥ 2.5 ¥ 103 5 ¥ 103
= 1.8 ¥ 1011 C / kg
E xamplE 6 An electron moves in the earth’s magnetic field of 5 ¥ 10–5 T with the energy of 10 keV. Find the
Larmour radius of the electron neglecting its velocity component parallel to the magnetic field.
Solution
1 2
E= mv^0 fi v^0 = 2 E / m
2
v mv
Formula used rL = ^0 = ^0
wc eB
2 ¥ 10 ¥ 103 ¥ 1.6 ¥ 10 -19

Now, v^ 0 = = 5.93 ¥ 107 m/s
9.1 ¥ 10 -31
9.1 ¥ 10-31 ¥ 5.93 ¥ 107

\ rL = = 6.75 m
1.6 ¥ 10-19 ¥ 5 ¥ 10-5
E xamplE 7 A solar wind proton is streaming with velocity v^0 = 3 ¥ 105 m/sec in the magnetic field of
5 ¥ 10 T. Compute the Larmour radius. Given that mass of proton = mp = 1.67 ¥ 10–27 kg.
–9
Solution Given B = 5 ¥ 109 T, mp = 1.67 ¥ 10–27 kg and v^0 = 3 ¥ 105 m/sec.

Larmour radius
m p v^ 0 1.67 ¥ 10-27 ¥ 3 ¥ 105
rL = =
eB 1.6 ¥ 10-19 ¥ 5 ¥ 10-9
= 6.26 ¥ 105 m
E xamplE 8 Calculate the area traced by the trajectory of a 1 keV He+ ion in the solar atmosphere near a
sunspot, where B = 5 ¥ 10–2 T.
Solution Energy E = 1 keV = 1 ¥ 103 ¥ 1.6 ¥ 10–19 J
1 2 2E
E= mv^0 fi v^0 =
2 mHe+
2 ¥ 103 ¥ 1.6 ¥ 10-19

fi v^0 = = 2.19 ¥ 105 m/sec.
4 ¥ 1.67 ¥ 10-27
Here we have used the mass of He+ as 4 mp.
v^0 mHe+
Now He+ shall have gyratory motion, whose Larmour radius rL =
eB
2.19 ¥ 105 ¥ 4 ¥ 1.67 ¥ 10-27
fi rL = = 0.183 m
1.6 ¥ 10-19 ¥ 5 ¥ 10-2
Since the trajectory of He+ is simply the circular motion, the area traced by this
= p rL2
= 3.14 ¥ 0.1832 = 0.105 m2
E xamplE 9 An electron is moving in uniform electric and magnetic fields which are perpendicular to each
other. Find the drift of the guiding centre if the magnitudes of the electric field and magnetic field be 100 V/m
and 10–3 T, respectively.
Solution Under the actions of uniform electric and magnetic fields, the trajectory of the electrons is slanted helix with
changing pitch. The drift of the guiding centre is
  
vB = ( E ¥ B )/ B 2
 E
(vB ) =
B
100
= -3 = 1 ¥ 105 m /sec
10
E xamplE 10 An ion engine has a 1 T magnetic field, and a collection of H+ (behaving collectively) is to
 
be shot out at an E ¥ B velocity of 1 ¥ 106 m/s. How much internal electric field must be present in the
collection of ions (H+)?
 
   E ¥ B
Solution E ¥ B drift is the drift of guiding centre; given by vE =
B2
E
vE = or E = BvE
B
or E = 1 ¥ 1 ¥ 106 = 106 V/m
E xamplE 11 Magnification related studies were conducted using Scanning Electron Microscope. By keeping
length of the scan on the Cathode Ray Tube fixed, the length of the scan on the specimen was halved. What
will be the ratio of new magnification to the old magnification?
Electron Optics 231
Solution Magnification M is defined as

Length of scan on CRT (C)
M=
Length of scan on specimen
C
M1 =
L1
C C 2C
Similarly, M 2 = = =
L2 L1/ 2 L1
M 2 2C L1
or = ¥ =2
M1 L1 C
E xamplE 12 Calculate the ratio of new focal length to the old focal length, if the number of turns in the coil
used to form magnetic lens in magnetic focusing is increased by 10% and other parameters are kept fixed.
N 11N
Solution New no. of turns = N + =
10 10
Focal length of magnetic lens
C1V
f =
N 2I 2
CV C1V
fi f1 = 21 2 and f2 = 2
N I Ê 11ˆ 2 2
ÁË ˜¯ N I
10
2
f 2 Ê 10 ˆ
fi = Á ˜ = 0.826
f1 Ë 11¯
It means the new focal length is reduced by increasing the number of turns of the coil.
Q.1 Specific charge of an electron is equal to

(a) –1.6 ¥ 10–19 C (b) 9.1 ¥ 10–31 kg (c) 1.76 ¥ 1012 C/kg (d) 1.6 ¥ 10–19C.
 
Q.2 In Thomson’s method for the determination of specific charge of an electron, the E and B fields are
applied such that
(a) electric force is double of magnetic force
(b) electron gets affected only by Lorentz force
(c) electric force and Lorentz force are in balance
(d) electron follows a helical path.
Q.3 The trajectory of an electron under the influence of a uniform magnetic field, when it is injected in the
perpendicular direction to the magnetic field is
(a) circular (b) helix (c) linear (d) none of these.
Q.4 Larmour radius rL of the electron motion in uniform and constant electric and magnetic fields depends
on
(a) mass and charge of the electron only (b) magnetic field only
(c) perpendicular velocity only (d) all of these.
Q.5 Electron optics is the branch of physics which is concerned with


(a) beam of electrons and their deflections by E field
 
(b) beam of electrons and their deflections by E and B fields
(c) the radiation when electron jumps from the excited state to ground state
(d) none of the above.
Q.6 Electrostatic focusing is achieved when
(a) constant electric field is applied to focus the electron beam
(b) electric field and magnetic field both are applied in the same direction
(c) two opposite electric fields and a non-uniform magnetic field are applied
Q.7 An example of magnetostatic focusing is
(a) electron microscope (b) electron gun
(c) anode (d) cathode.
Q.8 Lenses in electrostatic focusing are related to
(a) the electrostatic field areas (b) the magnetostatic field areas
(c) plane glass plates (d) convex lens used to focus the light.
Q.9 Magnetic lens formed for magnetostatic focusing suffer from defects like
(a) spherical aberration (b) chromatic aberration
(c) astigmatism (d) none of the above.
Q.10 In STM, the tunneling current decays
(a) linearly with the increase of gap between tip and sample
(b) exponentially with the increase of gap between tip and sample
(c) quadratically with the gap between tip and sample
P ractice P roblems
Q.1 Define electron optics and discuss the physics behind the coupling of the two words ‘electron’ and
‘optics’.
Q.2 Discuss J.J. Thomson’s method for the determination of specific charge of an electron.
Q.3 Derive the expression for cyclotron frequency of an electron under the action of a uniform magnetic
field. Discuss its trajectory based on proper mathematical expressions.
 
Q.4 Prove that the motion of an election in uniform E and B fields is the sum of simple gyratory motion
 
and the motion of guiding centre when both E and B fields are perpendicular to each other.
Q.5 Discuss electrostatic focusing of electron beam. Give an example of this focusing.
Q.6 How do you achieve magnetostatic focusing with the help of magnetic field. Discuss the similarities
between magnetic lens and the optical lens.
Q.7 Discuss the principle of SEM and the image formation.
Q.8 Discuss in short the applications of the SEM.
Q.9 Discuss the principle of STM and the behavior of tunneling current I with the gap between the tip and
the sample.
Waves and Oscillations 7
Learning Objectives LO 3 Know about simple pendulum, mass
string system, and damped harmonic
After reading this chapter you will be able to oscillator
LO 1 Differentiate between translational and LO 4 Explain attenuation coefficients of a
oscillatory motion vibrating system
LO 2 Learn about Simple Harmonic Motion LO 5 Discuss forced vibration
(SHM) and its differential equation LO 6 Understand resonance
Introduction
The motion of things can be broadly classified into two classes. It is according to whether the thing that
is moving stays near one place or travels from one place to another. The examples of things that stay
near one place are an oscillating pendulum, a vibrating violin string, electron vibration in atoms, etc. The
examples of things that travel from one place to another are a sliding hockey puck, a pulse traveling down
a long stretched rope plucked at one end, ocean waves rolling towards the beach, electron beam of a
television tube, etc. The motion of physical bodies may be classified mainly into two categories, namely
translational motion and vibrational or oscillatory motion.
7.1 TransLaTiOnaL MOTiOn LO1

If the position of a body varies linearly with time, for example, when a ball rolls on the ground or train moves
on a straight track, them such motions are called translational motions.
7.2 VibraTiOnaL Or OsciLLaTOry MOTiOn LO1

A motion of a body that repeats itself after regular intervals of times and when the body moves back and forth
over the same path is called vibrational or oscillatory motion. For example, the oscillations of the arms of a
walking person, the bob of the pendulum clock, beating of heart, etc. are in vibrational motion.
7.3 Simple Harmonic motion (SHm) LO2

If the acceleration of a particle in a periodic motion is always directly proportional to its displacement from
its equilibrium position and is always directed towards equilibrium position, then the motion of the particle
is said to be Simple Harmonic Motion (SHM).
7.3.1 type of SHm
Simple harmonic motion can be broadly classified into two classes, namely linear simple harmonic motion
and angular simple harmonic motion.
Linear Simple Harmonic Motion
The motion is said to be linear simple harmonic motion, if the displacement of a particle executing SHM is
linear. The examples are the motion of simple pendulum, the motion of a point mass tied with a spring, etc.
Angular Simple Harmonic Motion
The motion is said to be angular simple harmonic motion, if the displacement of a particle executing SHM is
angular. The examples of angular SHM are torsional oscillations and oscillations of a compound pendulum.
7.3.2 essential condition of SHm

If f be the linear acceleration and x the displacement from the equilibrium position, then essential condition
for linear SHM is
f µ –x
If a be the angular acceleration and q the angular displacement from the equilibrium position, then the
condition of angular SHM is
aµ–q
Time Period: The smallest time interval during which the oscillation repeats itself is called the time period.
It is denoted by T and its unit is seconds.
Frequency: The number of oscillations that a body completes in one second is called the frequency of
periodic motion. It is the reciprocal of the time period T and is given by
1
n=
T
unit of frequency is per second or Hertz, represented by Hz.
Amplitude: The amplitude of an SHM is the maximum displacement of the body from its mean position.
Phase: It is the physical quantity that expresses the instantaneous position and direction of motion of an
oscillating system. The quantity (wt + d) of the sine function is called the phase of the motion, if the motion
is represented by y = A sin (wt + d), together with w as the angular frequency.
7.4 DifferenTiaL equaTiOn Of sHM anD iTs sOLuTiOn LO2

Let us consider a particle of mass m executing an SHM along a straight line with x as its displacement from
the mean position at any time t. Then from the basic condition of the SHM, the restoring force F will be
proportional to the displacement and will be directed opposite to this. Therefore,
Fµ–x or F = –kx (i)

where k is a proportionality constant and is known as force constant.
d 2x
If f = be the acceleration at any instant of time, then
dt 2
d 2x
m 2 = - kx
dt
2
d x k
+ x=0 (ii)
dt 2 m
k
Putting = w2, we have
m
d 2x
+ w2x = 0 (iii)
dt 2
The above Eqs. (ii) and (iii) are known as the differential equations of the SHM
In order to get the solution of Eq. (iii), we assume x as
x = Cea t (iv)
where C and a are constants. Then
dx
= Ca ea t
dt
d 2x
and = Ca 2ea t (v)
dt 2
By using Eqs. (iv) and (v), Eq. (iii) becomes
Ca2ea t + w2Ceat = 0
Ceat (a2 + w2) = 0
Ceat π 0, as it is the displacement x
a2 + w2 = 0
or a = ±iw
where i is the imaginary number, given by i = -1 .
Since a has two values, possible solutions of Eq. (iii) are
x = Ceiwt and x = C–iwt,
the linear combination of which will give the general solution Eq. (iii) as
x = C1reiwt + C2e–iwt (vi)
where C1 and C2 are constants.
Equation (vi) can also be written as
x = C1(cos wt + i sin wt) + C2(cos wt – i sin wt)
= (C1 + C2) cos wt + (iC1 – iC2) sin wt (vii)
Let C1 + C2 = A sin d and iC1 – iC2 = A cos d

Then Eq. (vii) gives
x = A sin d cos wt + A cos d sin w t
x = Asin (wt + d) (viii)
This is the desired solution of Eq. (iii), which gives the displacement of the particle executing SHM at any
instant of time. Here A is the maximum displacement of the particle
È 2p ˘
In Eq. (viii), we replaced t by Ít + , then we have
Î w ˙˚
È Ê 2p ˆ ˘
x = A sin Íw Á t + ˜ +d˙
Î Ë w¯ ˚
= A sin (wt + 2p + d)
or x = A sin (wt + d)
Here, it may be noted that this is the same as Eq. (viii) and shows that the motion is repeated after an interval
2p 2p
of . So this interval will be the time period of the SHM, given by T =
w w
As we know that
k k
w2 = fiw =
m m
2p m
\ T= = 2p (ix)
w k
Frequency:
1 w 1 k
n= = = (x)
T 2p 2p m
Phase: The quantity (wt + d) is known as the phase of the vibrating particle. If t = 0 then wt + d = d, so that
initial phase will be d. If a particle starts motion from its mean position them d will be zero but if it starts
motion from the extreme position then d will be p/2.
Velocity and acceleration: The particle executing SHM is a harmonic oscilliator.
We can find its velocity from the expression of its displacement, given below
x = A sin (wt + d)
Differentiating it w.r.t. time, we get
dx
v= = Aw cos (wt + d) (xi)
dt
= Aw 1 - sin 2 (w t + d ) = w A2 - A2 sin 2 (w t + d )
2 2
or v= w A -x
This is the expression for velocity of the particle at any displacement x. The maximum velocity is obtained
by putting x = 0.
\ vmax = wA
Since x = 0 corresponds to its mean position, the particle has maximum velocity when it is at the mean
position. At the maximum displacement, i.e., at the extreme position of the particle the velocity is obtained
as zero. This extreme position is
x=A
Differentiating Eq. (xi) w.r.t. time t , we get
dv
f= = –Aw2 sin (wt + d)
dt
or f = –w2x (xiii)
The above equation gives acceleration of the oscillating particle at any displacement. This equation is the
standard equation of SHM
This is clear from Eq. (xiii) that for the maximum acceleration
x = A (the extreme position)
\ Maximum acceleration,
fmax = w2A at the extreme position
Minimum acceleration is obtained by putting x = 0
fmin = 0 at the mean position
7.4.1 total energy of a Harmonic oscillator

The total energy of a simple harmonic oscillator at any point would be the sum of kinetic energy (K.E.) and
potential energy (P.E.), as it executes vibrations under the action of a free restoring force.
The potential energy of a system is given by the amount of work done to move the system from the position
0 to x by applying a force. Therefore,
x
  1
P.E. = Ú F ◊ dx Ú kx dx = 2 kx
2
(xiv)
0
Here, we have not put the negative sign, as it simply tells us whether the work is done by the system or on the
system. Putting the values of k and x, we get
1
P.E. = mw2 A2 sin2 (wt + d) (xv)
2
The kinetic energy of the harmonic oscillator is
2
1 2 1 È dx ˘
K.E. = mv = m Í ˙
2 2 Î dt ˚
1
= mw 2 A2 cos 2 (w t + d ) (xvi)
2
The total energy E of the oscillator at any displacement can be obtained as

E = P.E. + K.E.
1 1
= mw2 A2 sin2(wt + d) + mw2 A2 cos2 (wt + d)
2 2
1
E= mw2 A2 (xvii)
2
Since the total energy E is constant, it is obvious that the maximum possible value of K.E. or P.E. would be the
1
same as mw2 A2. It is interesting to learn how the kinetic and potential energies of the harmonic oscillator
2
vary with time, if the oscillations start at t = 0 from its extreme position. This is shown in Fig. 7.1.
Potential Energy
E E = (1/2)mw 2A2
Kinetic Energy
O t
T/2 T
Figure 7.1
The average potential energy of the simple harmonic oscillator for one complete cycle is
T
1 1 2
T Ú0 2
<P.E.> = kx dt
T
1
Ú 2 kA
2
sin 2 (w t + d ) dt
= 0
T
T
1
2 Ú0
mw 2 A2 sin 2 (w t + d ) dt
=
T
T
Ú sin
2
(w t + d ) dt
1
= mw 2 A2 0
2 T
È T 2 ˘
1 Í Ú sin (w t + d ) dt ˙
<P.E.> = mw 2 A2 Í 0 1˙
4 = ˙
ÍÎas T 2˚
Also, the average kinetic energy for one complete cycle is

T 2
1 1 Ê dx ˆ
T Ú0 2 Ë dt ¯
<K.E.> = m Á ˜ dt
T
1
Ú 2 mA w
2 2
cos 2 (w t + d ) dt
= 0
T
T
Ú cos
2
(w t + d ) dt
1
= mw 2 A2 0
2 T
1
<K.E.> = mw 2 A2
4
1
fi <K.E.> = <P.E.> = mw 2 A2
4
From the above calculations, it is clear that the average kinetic energy is equal to the average potential energy
for a harmonic oscillator over a complete cycle
Now the total average energy over a complete cycle is
1 1
<E> = <K.E.> + <P.E.> = mw2 A2 + mw2 A2
4 4
1
mw2 A2
=
2
which is equal to the total energy of the harmonic oscillator.
7.5 siMpLe penDuLuM LO3

An ideal simple pendulum is formed when a heavy point mass is A
suspended, by a weightless, inextensible and perfectly flexible string
from a rigid support. However, these requirements cannot be fulfilled a
in practice. Therefore, a simple pendulum consists of small heavy T
metallic sphere suspended by a long fine string from a rigid support.
l
The arrangement of a simple pendulum is shown in Fig. 7.2, where P
pendulum AO is in equilibrium position. When the mass of the
pendulum is displaced from its equilibrium position and released, it
starts oscillating around its equilibrium position. The motion of the
mg sin a mg cos a
point mass is said to be SHM Suppose at any instant the particle is
O mg
at point P at a distance x from its equilibrium position and makes an
small angle a. The length of the string is l. Figure 7.2
At point P, the particle is under the action of various forces: (i) the
weight mg of the bob acting vertically downward (ii) tension T in string along PA. The weight mg has two
components (a) radial component mg cos a, (b) transverse component mg sin a.
It is clear from the above arrangement that the tension in the string is opposed by the radial component mg cos
a. Therefore, the force T – mg cos a provides centripetal force for circular arc and the tangential component
mg sin a tends to bring the bob back to its initial position. Thus mg sin a is often known as restoring force
and therefore
F = –mg sin a (i)
2
d x
The negative sign indicates that the acceleration and the displacement are oppositely directed. If be the
dt 2
acceleration at any time t in the direction of increasing x, then the force
d 2x
F= m (ii)
dt 2
Therefore,
d 2x
m = - mg sin a
dt 2
d 2x
or = - g sin a (iii)
dt 2
For small angle a, the distance x (arc) can be written in terms of l and a, as
x = la
which on differentiation gives
d 2x d 2a
= l
dt 2 dt 2
d 2a
or l = - g sin a
dt 2
d 2a g
or = - sin a
dt 2 l
d 2a g
or + sin a = 0 (iv)
dt 2 l
Now, we realize that Eq. (iv) is the equation of motion of the pendulum. For a small deflection a, we can
write sin a ª a. Then
d 2a g
+ a =0
dt 2 l
The above equation is known as equation of motion of simple pendulum, whose solution can be written as
a = a0 sin (wt + d)
1/2
Ê gˆ
where, w = Á ˜ and d is the initial phase.
Ël¯
Therefore, the time period is given by
2p l
T= = 2p
w g
From the above expression, it is clear that the time period of a simple pendulum is independent of mass and
shape of the bob.
7.6 maSS-String SyStem LO3

When a spring is compressed or stretched by applying a force on it, then the resulting compression does not
have a simple linear relationship with the applied force. The relationship is linear only for small displacements.
Here, the elastic force F produced in the spring is given by
F = –kx
together with x as the change in length of the spring and k as the spring constant. Now we shall discuss two
simple cases.
7.6.1 Horizontal oscillations
Here we assume a massless spring, one end of which is connected to a mass m and the other end is connected
to a fixed point, as shown in Fig. 7.3
Relaxed
F=0
(a)
F
Stretched
(b)
F x
Compressed
(c)
x
Figure 7.3
The mass m is free to move on a frictionless horizontal surface. The static equilibrium position is shown in
Fig. 7.3(a) as relaxed and no force is acting on it. When the mass m is pulled to the right [Fig. 7.3(b)], through
a small distance x, then the restoring force exerted by the spring is directed towards the left and is given by
F = – kx
where the negative sign indicate that the force and displacement are oppositely directed. Here, the mass starts
d 2x
moving with linear acceleration 2 . Then, we have
dt
d 2x
F =m 2
dt
d 2x
or - kx = m 2 [By Newton’s second law]
dt
d 2x k d 2x
or + x = 0, i.e., + w 2x = 0
dt 2 m dt 2
This is the equation of SHM, having time period

2p m
T= = 2p
w k
7.6.2 Vertical oscillations
Let us consider a perfectly elastic and massless spring of length L hanging freely from a support, as shown in
Fig. 7.4(a). When a mass m is attached to its lower end, it is stretched through a distance x ¢ by the force mg.
F = – kx ¢
where k is the force constant of the spring. The another force mg (weight) is acting downward on the spring.
Since in this situation no net force acts on the body of mass m, i.e., mg = –kx ¢. Now we pull down the mass m
through a small distance y ¢ from the equilibrium position and release, then its starts oscillating, as shown in
Fig. 7.4(c), (d). F = –ky ¢ is the restoring force and is oppositely directed to the displacement.
Static Equilibrium
Loaded
Relaxed
¢ ¢
¢
F
(a) (b) (c) (d)
Figure 7.4
By applying Newton’s second law, we get

d 2 y¢
F =m = - ky ¢
dt 2
d 2 y¢ k
or 2
= - y¢ = - w 2 y¢
dt m
k
where w2 = . This is the equation of motion of mass m executing SHM with time period
m
2p m
T= = 2p
w k
By using this formula, we can calculate the time period of mass-spring system.
7.7 DaMpeD HarMOnic OsciLLaTOr LO3

In the previous section free oscillations were discussed in some oscillating systems. In
such systems, no frictional force or resistance is offered. Therefore, the body will keep on
vibrating indefinitely and such vibrations are called free vibrations. But in real situation, there
is always some resistance offered to the oscillating system. In real sense if a body set into
vibrations will have its amplitude continuously decreasing due to fractional resistance and so
the vibrations will die after some time. The motion is said to be damped by the friction and is
m
called as damped vibrations. For example, when a pendulum is displaced from its equilibrium
position, it oscillates with a decreasing amplitude and finally comes to the rest. Figure 7.5
Suppose a system has a body of mass m attached to a spring whose force constant is k
(Fig. 7.5). Again consider x ¢ to be the displacement of the body from the equilibrium state at any instant of
dx ¢
time. Then will be the instantaneous velocity.
dt
Here two types of forces are acting on the body: a restoring force (Fr = –kx ¢) which is proportional to the
dx ¢
displacement x ¢ and acts in the opposite direction to the displacement and a damping force - q ¢ which
dt
is proportional to the velocity and is oppositely directed to the motion. The total force acting on the body is
thus given by
dx ¢
F = - kx ¢ - q ¢
dt
d 2 x¢ dx ¢
or m 2
= - kx ¢ - q ¢
dt dt
d 2 x ¢ q ¢ dx ¢ k
or + + x¢ = 0 (i)
dt 2 m dt m
q¢ k
If we put = 2s and = w2, then above equation takes the form
m m
d 2 x¢ dx ¢
2
+ 2s + w 2 x¢ = 0 (ii)
dt dt
which is the equation of second degree. We assume its solution as
x ¢ = Aeat (iii)
where A and a are arbitrary constants.
Differentiating Eq. (iii) w.r.t. t, we have
dx ¢ d 2 x¢
= Aa ea t and 2 = Aa 2ea t
dt dt
By putting these values in Eq. (ii), we have
Aa2eat + 2sAaeat + w2Aeat = 0
or Aeat [a2 + 2sa + w2] = 0
 Aeat π 0
\ a2 + 2sa + w2 = 0
This gives,
a = - s ± s2 - w 2
Therefore, a has two roots
a1 = - s + s 2 - w 2
a 2 = - s - s2 - w 2
Therefore, the general solution of Eq. (ii) is given by

s 2 - w 2 )t s 2 - w 2 )t (iv)
x ¢ = A1e( - s + + A2e( - s -
where A1 and A2 are arbitrary constants. The actual solution depends upon whether s2 > w2, s2 = w2 and
s2 < w2. Now we will discuss all these cases.
Case-A: s2 > w2
2 2
The term s 2 - w 2 is real and less than s. Therefore both exponents i.e., [- s + s 2 - w 2 ] and [- s - s - w ]
in Eq. (iv) are negative. Due to this reason the displacement x ¢ continuously decreases exponentially to zero
without performing any oscillation [Fig. 7.6(a)]. This motion is known as overdamped.
Case-B: s2 = w2
For this case, Eq. (iv) does not satisfy differential Eq. (ii).Suppose that s 2 - w 2 is not exactly zero but it is
equal to very small quantity b. Now, Eq. (iv) gives
x ¢ = A1 exp (–s + b)t + A2 exp (–s – b) t
x ¢ = e–st[A1ebt + A2e–bt]
= e–st [A1(1 + bt + …) + A2(1 – bt + …)]
Here we have neglected the terms containing b2, b3, b4 is very small
\ x ¢ = e–st[(A1 + A2) + bt(A1 – A2) + …]
x ¢ = e–st[P ¢ + Q ¢t]
where P ¢ = (A1 + A2) and Q ¢ = b(A1 – A2) (v)
Equation (v) represents a possible solution of this form. From this equation, it is clear that as t increases
the term (P ¢ + Q ¢t) increases but the term e–st gets decreased. Because of this fact the displacement x ¢
first increases due to the term (P ¢ + Q¢t) but it decreases because of the exponential term e–st and finally it
approaches to zero as t increases. When we compare cases A and B it is noticed that the exponent is –st in case
B while it as more than –st in case of A. Therefore, in case A the particle acquires its position of equilibrium
rapidly than in case B[ Fig. 7.6(b)]. This type of motion is called critical damped motion. Examples are
voltmeter, ammeter, etc. In these instruments, pointer moves to the correct position and comes to the rest
without any oscillations.
Case-C: s2 < w2
The term s 2 - w 2 is imaginary, which can be written as
s 2 - w 2 = i w 2 - s 2 = ib ¢
where b ¢ = w 2 - s 2 and i -1
Now, Eq. (iv) becomes
x ¢ = A1e(–s + ib ¢ )t + A2e(–s – ib ¢ )t
x ¢ = e–st[A1eib ¢t + A2e–ib ¢t]
= e–st [A1 cos b ¢t + i sin b ¢t) + A2(cos b ¢t – i sin b ¢t)
= e–st[(A1 + A2) cos b ¢t + i[(A1 – A2) sin b ¢t]
= e–st[(A sin d cos b ¢t + A cos d sin b ¢t]
where A sin d = A1 + A2 and A cos d = i(A1 – A2)
\ x¢ = e –st
A sin (b ¢t + d)
Putting the value of b ¢ in the above questions, we get
x ¢ = Ae–st sin [ (w 2 - s 2 )t + d ] (vi)

The above equation shows the oscillatory motion and represents the damped harmonic oscillator. The
oscillations are not simple harmonic because the amplitude (Ae–st) is not constant but decreases with time (t).
However, the decay of amplitude depends upon the damping factor s. The motion is known as under damped
motion [Fig. 7.6(c)]. The motion of a pendulum in air, the motion of the coil of ballistic galvanometer, etc.
are the example of this case.
(b)
(a)
(c)
Figure 7.6
7.8 attenuation coefficientS of a Vibrating SyStem LO4

By the definition, the energy of an oscillator is proportional to the square of its amplitude. In damped
oscillator, the amplitude decays exponentially with time as e–st. It means that the energy will decay as (e–st)2,
i.e., e–2st. So the decay rate of energy depends upon s. The following three characteristics namely logarithmic
decrement, relaxation time and quality factor may give the attenuation of a vibrating system.
7.8.1 logarithmic Decrements
The rate at which the amplitude dies away is measured by logarithmic decrement. The amplitude of damped
harmonic oscillator is given by a factor Ae–st. Therefore, at time t = 0 the amplitude will be maximum (i.e.,
A = A0). If A1, A2, A3, … be the amplitude at time t = T, 2T, 3T, … respectively where T is the time period of
oscillations, then.
A1 = Ae–sT, A2 = Ae–s(2T), A3 = Ae–s(3T) and so on.
This yields
A0 A1 A2
= = = … = esT = el (where sT = l)
A1 A2 A3
Here, l is called logarithmic decrement.
Now, by taking the natural logarithmic, we have
A0 A A
ln = ln 1 = ln 2 = l
A1 A2 A3
Hence, logarithmic decrement is the natural logarithm of ratio between two successive maximum amplitudes,
which are separated by one period.
7.8.2 relaxation time
1
It is the time taken by damped harmonic oscillator for decaying total mechanical energy by the factor of
e
its initial value.
The mechanical energy of the oscillator is
1
E= mA2w 2e -2 st (i)
2
1
At t = 0, E = E0 = mA2w 02
2
\ Total energy, E = E0e–2st (ii)
Suppose t be the relaxation time, then at time t = t,
E
E= 0 (By definition)
e
By using Eq. (ii), we get
E0
= E0e -2 st
e
or el = e2st
or 1 = 2st
1
or t= (iii)
2s
Therefore, the dissipated energy in terms of relaxation time is written as

E = E0e–t/t (iv)
7.8.3 Quality factor
It is defined as 2p times the ratio of energy stored in the system to the energy lost per cycle. This factor of a
damped oscillator shows the quality of oscillator so far as damping is concerned.
E
Q = 2p
Pd T
(v)
where Pd is the power dissipation and T is the periodic time. Then,
E 2pt È E˘
Q = 2p = ÍÎ Pd = t ˙˚
( E /t )T T
È 2p ˘
Q = wt ÍÎ w = T ˙˚ (vi)
From the above equation, it is clear that the value of relaxation time t will be higher (or damping will be
lower) for higher value of Q.
For the force constant k and the mass m of the vibrating system
k 1
w= and t = [from Eq. (iii)]
m 2s
1 k
\ Q=
2s m
Since lower values of s lead to lower damping, it is clear that for low damping, the quality factor would be
higher.
7.9 fOrceD VibraTiOns LO5

We have discussed earlier the vibrations in which the body vibrates with its own frequency without being
placed to any external force. The different situation arises if the body is placed to an external force while it
is vibrating. These vibrations are known as forced vibrations. For example, if a bob of simple pendulum is
held in hand and then given number of swings by the hand. In this case, the pendulum vibrates due to external
force and not due to its natural frequency. So forced vibrations can also be defined as the vibrations in which
the body vibrates with frequency other than its natural frequency, which is due to some external periodic
force. The tuning fork is also one of the examples of forced vibrations.
Theory of Forced Vibrations
Suppose a particle of mass m is connected to a spring. When it is displaced from its mean position, the
oscillations are started and the particle experiences different kinds of forces viz, a restoring force (–kx), a
Ê dx ˆ
damping force Á - q ¢ ˜ and the external periodic force (F0 sin wt). The total force acting on the particle is,
Ë dt ¯
therefore
dx
F = F0 sin wt – q ¢ - kx (i)
dt
By Newton’s second law of motion
d 2x
F =m k
dt 2 F0 sin w t
Hence,
d 2x dx
m = F0 sin w t - q ¢ - kx
dt 2 dt Figure 7.7
2
d x q ¢ dx k F sin w t
or 2
+ + x= 0 (ii)
dt m dt m m
Eq (ii) is the differential equation of motion of the particle.
q¢ k F
Substitute = 2 s, = w 02 and 0 = f , then Eq, (ii) becomes
m m m
d 2x dx
2
+ 2 s + w 02 x = f sin w t (iii)
dt dt
In the steady state, the solution of the above equation should be
x = A sin (wt – d) (iv)
where A is the amplitude of vibrations in the steady state. By differentiating Eq. (iv) twice w.r.t. t, we have
dx
= w A cos (w t - d )
dt
d 2x
and = - w 2 A sin (w t - d )
dt 2
dx d 2x
By substituting the values of x, and 2 in Eq. (iii), we have
dt dt
– w2 A sin (wt – d) + 2s w A cos (wt – d) + w 02 A sin (wt – d)
= f sin {(wt – d) + d}
or A( w 02 – w2) sin (wt – d) + 2sw A cos (wt – d)
= f sin (wt – d) cos d + f cos (wt – d) sin d (v)
If Eq. (v) holds good for all values of t, then the coefficients of sin (wt – d) and cos (wt – d) must be equal
on both the sides, then
A( w 02 – w2) = f cos d (vi)
and 2s wA = f sin d (vii)
By squaring and adding Eqs. (vi) and (vii), we have
2
A2( w 0 – w2)2 + 4s2 w2A2 = f2
f
or A= (viii)
(w 02 - w ) + 4 s 2w 2
2 2
Dividing Eq. (vii) by Eq. (vi), then we get

2 sw
tan d =
w 02 - w 2
È 2 sw ˘
or phase, d = tan -1 Í 2 2˙
(ix)
Îw0 - w ˚
From Eq. (viii) and (ix), it is clear that the amplitude and phase of the forced oscillations depend upon
( w 02 – w2), i.e., these depend upon the driving frequency (w) and the natural frequency (w0) of the oscillator.
The amplitude and the phase are explained as below.
Case - A: Very low driving frequency, i.e., w < < w0. In this case, the amplitude of the vibrations
f f
A= ª
(w 02 - w 2 ) 2 + 4 s 2w 2 w 02
F0 È F0 ˘
= ÍÎ f = m ˙˚
mw 02
F0 È k˘
or A= 2
ÍÎ w 0 = m ˙˚
k
Hence, the amplitude depends on the force constant of the spring and the magnitude of the applied force.
È 2 sw ˘ -1 È 2 sw ˘
Phase, d = tan -1 Í 2 2 ˙ = tan Í 2 ˙
Îw0 - w ˚ Î w0 ˚
Since w 02 > > w, 2sw/w02 Æ 0 and d Æ 0 or ª 0. Therefore, under this situation the driving force and the
displacement are in phase.
Case - B: Same driving and natural frequencies, i.e., w = w0. This frequency is called resonant frequency.
Under this situation, the amplitude of vibrations
f
A=
(w 02 - w 2 ) 2 + 4s 2w 2
f F0 /m F
= = = 0
2 sw 0 (q ¢ / m)w 0 q ¢w 0
Hence, the amplitude of vibrations depends upon the damping and applied force. Now
È 2 sw ˘ -1 È 2 sw ˘
Phase, d = tan -1 Í 2 2 ˙ = tan Í
(w
Î 0 - w ) ˚ Î 0 ˙˚
p
= tan -1[•] =
2
p
Thus, the displacement lags behind the force by a phase of , as x = A sin (wt – d) and the applied force if
2
F0 sin wt.
Case - C: Very large driving frequency, i.e., w > > w0. Here, the amplitude of vibrations
f
A=
(w 02 - w ) + 4 s 2w 2
2 2
Since w is very large w4 > > 4w2s2

f F
A= = 0
w 2 mw 2
È 2 sw ˘ -1 È 2 sw ˘
Phase, d = tan -1 Í 2 2 ˙ = tan Í
(w
Î 0 - w ) ˚ Î -w 2 ˚˙
È 2s ˘
ª tan–1 Í = tan–1[–0] = p [Since w is very large, 1/w = 0]
Î -w ˙˚
Therefore, under the situation w > > w0, the displacement lags behind the force by a phase of p.
7.10 resOnance LO6

So far we have discussed various types of oscillations. It is clear that if a vibrating system consists of a
number of oscillators coupled together, the resultant motion would be complicated. However, if we choose
the starting conditions correctly, it is possible to cause the system vibrate in such a way that every part has the
same frequency. Such simple vibrations are known as normal oscillations or normal modes and the associated
frequencies are called the normal frequencies or the natural frequencies. So this is evident that all mechanical
structures, for example, buildings, airplanes, bridges, etc. have one or more natural frequencies. Now if such
a structure is subject to a driving frequency, which is equal to one of the natural frequencies, the resulting
oscillations will have large amplitude that can have disastrous consequences. Shattering a wine glass with a
sound wave that matches one of the natural frequencies of the glass is one demonstration of this phenomenon
of resonance. Another outcome of this effect is the collapse of roadways and bridges in earthquakes.
In order to make clear the phenomenon of resonance, we take an example of forced oscillations that occur at
the frequency of the external force (driving frequency w) and not at the natural frequency (w0) of the vibrating
system. As has been seen, however, the amplitude of oscillations depends on the relationship between the
driving frequency of the applied force and the natural frequency. A succession of small impulses, if applied
at the proper frequency, can develop an oscillation of large amplitude. You can take an example of pushing a
friend on a swing. You would have noticed that by applying pushes precisely at the same time in each cycle,
you can cause your friend to move in an increasingly large arc.
In Fig.7.8, we show the variation of amplitude of the forced vibrations with the driving frequency w for
three cases of large damping, medium damping and small damping. Medium damping corresponds to twice
the damping force and large damping corresponds to four times the damping force. It is clear from the
figure that the amplitude decreases if the frequency w is far from the condition w = w0, called condition of
Small Damping
Medium Damping
Amplitude
Large Damping
Natural Frequency
w 0/2 w0 3w 0/2 2w 0
Driving Frequency w
Figure 7.8
resonance. Moreover, when the damping is small, the amplitude of the forced oscillations increases rapidly
as w approaches w0. The amplitude reaches its maximum when w = w0. For medium damping also, the
amplitude gets increased but it does not increase so rapidly near the resonance (w = w0). However, for the
largest damping the resonant frequency is displaced slightly from the natural frequency.
s UmmarY
The topics covered in this chapter are summarised below.
✦ The motion of physical bodies is broadly classified into two categories, namely translational motion
and vibrational or oscillatory motion. If the position of a body varies linearly with time, then such
motions are called translational motions. The examples of translational motion are a ball that rolls on
the ground and a train that moves on a straight track. A motion of a body that repeats itself after regular
intervals of times and when the body moves back and forth over the same path is called vibrational
or oscillatory motion. The example of vibrational motion are the oscillations of the arms of a walking
person, the bob of the pendulum clock, beating of heart, etc.
✦ If the acceleration of a particle in a periodic motion is always directly proportional to its displacement
from its equilibrium position and is always directed towards equilibrium position, then the motion of
the particle is said to be Simple Harmonic Motion (SHM).
✦ For f as the linear acceleration of the particle and x as its displacement from the equilibrium position,
the essential condition for linear SHM is f µ – x. However, if a be the angular acceleration and q be the
angular displacement from the equilibrium position, then the condition for angular SHM is a µ – q.
✦ For the displacement x and the angular frequency w ( = k /m , where k is the force constant and m is
d 2x
the mass of the particle), 2 + w2x = 0 represents the differential equation of the SHM
dt
✦ The solution x = A sin (wt + d) gives the displacement of the particle executing SHM at any instant of
time t. Here A represents the maximum displacement of the particle, which is called the amplitude of
2 2
oscillations. The velocity of the particle is given by v = w ( A - x ) and the acceleration is f = –w2x.
1
The energy of a harmonic oscillator is given by E = mw2 A2.
2
✦ When there is no frictional force or resistance, the body will keep on vibration indefinitely and such
vibrations are called free vibrations. But in real situation, there is always some resistance offered to
the oscillating system. Then a body when sets into vibrations will have its amplitude continuously
decreasing due to fractional resistance and hence the vibrations will die after some time. The motion is
said to be damped by the friction and is called as damped vibrations.
✦ The energy of an oscillator is proportional to the square of its amplitude. In damped oscillator, the
amplitude decays exponentially with time as e–st, where s = q ¢/2m together with m as the mass of the
body and q ¢ as the proportionality constant of the damping force. Accordingly the energy also decays.
So the decay rate of energy depends upon s. In this context, the three characteristics namely logarithmic
decrement, relaxation time and quality factor give the attenuation of a vibrating system.
✦ If A0, A1, A2, A3, … be the amplitudes at time t = 0, T, 2T, 3T, …, respectively, where T is the time
A A A
period of oscillations, then logarithmic decrement is defined as l = ln 0 = ln 1 = ln 2 .
A1 A2 A3
✦ The relation time is the time taken by damped harmonic oscillator for decaying total mechanical energy
1
the factor of 1/e of its initial value. It is given by t = where s = q ¢/2m together with m as the mass
2s
of the body and q ¢ as the proportionality constant of the damping force.
✦ Quality factor Q of the oscillator is defined as 2p times the ratio of energy stored in the system to the
1 k
energy lost per cycle. This factor shows the quality of the oscillator and is given by Q = wt = .
2s m
The higher value of Q means lower damping of the oscillator.
✦ All mechanical structures, for example, buildings, airplanes, bridges, etc. have one or more natural
frequencies. If such a structure is subject to a driving frequency (say w), which is equal to one of the
natural frequencies (say w0), the resulting oscillations will have large amplitude that can have disastrous
consequences. Shattering a wine glass with a sound wave that matches one of the natural frequencies
of the glass is one demonostration of this phenomenon of resonance. Another outcome of this effect is
the collapse of roadways and bridges in earthquakes. The condition (for example, for forced oscillation)
where the driving frequency w and the natural frequency w0 of the vibrating system match with each
other is known as resonance. At the resonance, the amplitude of oscillations reaches its maximum.
s olVeD e XamPles
E xamplE 1 The total energy particle executing a SHM of period 2p seconds in 10.24 ¥ 10–4 Joule. The
displacement of a particle at p/4 second is 0.08 2m . Find the amplitude and mass of the particle.
Solution Given E = 10.24 ¥ 10–4 J, T = 2p sec and x = 0.08 2m at t = p/4 sec.
In SHM, the displacement of a particle is given by
x = A sin wt = A sin ÊÁ 2p t ˆ˜
Ë T ¯
2p p p A
0.08 2 = A sin = A sin =
2p 4 4 2
or A = 0.16 m
2mp 2 A2 ET 2
Total energy is given by E = or m =
T2 2p 2 A2
-4
10.24 ¥ 10 ¥ (2p ) 2
20.48 ¥ 10 -4
or m= 2 2
=
2p ¥ (0.16) 0.0256
= 0.08 kg
m = 80 g
E xamplE 2 A particle executes SHM of period 10 sec. and amplitude 5 cm. Calcualte the maximum
amplitude of velocity
Solution Given displacement amplitude (A) = 0.05 m and T = 10 s.
2p
Formula used for maximum amplitude of velocity = Aw = 0.05 ¥
T
0.05 ¥ 2 ¥ 3.14
= = 0.0314 m/s
10
E xamplE 3 Calculate the force constant and time period, if the potential of a harmonic oscillator of mass
2 kg in its resting position is 5.0 J and total energy is 9.0 J, when the amplitude is 1.0 m.
Solution Given E = 9.0 J, U = 5.0 J, K.E = E – U = 4.0 J and A = 1.0 m.
The kinetic energy of maximum displacement will be
1 1
K.E. = kA2 or 4.0 J = k(1.0)2
2 2
k = 8.0 J/m
m 2
T = 2p = 2p = p = 3.14 s
k 8
E xamplE 4 A particle is executing SHM of amplitude 0.06 m and a period of 6 s. Find out the time taken
by it in moving from one end of its path to a position 0.03 m from the equilibrium position on the same side.
Solution Given A = 0.06 m, T = 6.0 s and x = 0.03 m
Displacement of a particle executing simple harmonic motion
x = A sin (wt + d) (i)
At t = 0, particle is at one end so that at t = 0, x = A, then
By using Eq. (i)
A = A sin (0 + d) or sin d = 1 = sin p/2
So, d = p/2
Putting this value of d in Eq. (i), we get
Ê pˆ 2p
x = A sin ÁË w t + ˜¯ = A cos wt = A cos t
2 T
2p
or 0.03 = 0.06 cos t
6.0
Êp ˆ 1 Êpˆ
or cos Á t ˜ = = cos Á ˜
Ë3 ¯ 2 Ë 3¯
p p
or t=
3 3
or t = 1.0 s
E xamplE 5 Find the maximum velocity and acceleration of a particle executing SHM of period 10p second
and amplitude 5 ¥ 10–2 m.
Solution Given T = 10p sec and A = 5 ¥ 10–2 m.
The equation of simple harmonic motion is x = A sin (wt + d)
dx
v= = A w cos (wt + d)
dt
v will be maximum for cos (wt + d) = 1
2p 2p
Vmax = Aw = 5 ¥ 10–2 ¥ = 5 ¥ 10–2 ¥
T 10p
= 1.0 ¥ 10–2 m/s
d 2x
Acceleration (f) = = –Aw2 sin (wt + d)
dt 2
f will be maximum for sin (wt + d) = 1, then
2
Ê 2p ˆ
f = Aw2 = 5.0 ¥ 10–2 ¥ ÁË ˜¯
T
2
Ê 2p ˆ
= 5.0 ¥ 10–2 ¥ Á
Ë 10p ˜¯
= 2.0 ¥ 10–3 m/sec2
E xamplE 6 Calculate the maximum velocity of a particle that executes SHM of amplitude 0.06 m with time
period of 10 p s.
Solutioin Given A = 0.06 m and T = 10p s.
The equation of SHM is x = A sin (wt + d)
dx
v= = Aw cos (wt + d)
dt
v will b maximum if cos (wt + d) = 1
2p
vmax = Aw = 6.0 ¥ 10–2 ¥
10p
= 1.2 ¥ 10–2 m/s
E xamplE 7 A mass of 1.0 kg is attached to a spring of stiffness constant 16 N/m. Find the natural frequency
Solution Given k = 16 N/m and m = 1.0 kg.
1 k 1 16 2
Fornula used for natural frequency n = = =
2p m 2p 1 p
n = 0.64 Hz
E xamplE 8 A simple pendulum of one meter length is hanging at one end. Considering the oscillations to be
of small displacement, find the period of oscillation if the mass of pendulum is 2.0 kg.
Solution Given l = 1.0 m and m = 2.0 kg
l 1.0 1.0
Time period T = 2p = 2p = 2 ¥ 3.14 ¥
g 9.8 9.8
= 2.0 s
E xamplE 9 A particle of mass 100 gm is placed in a field of potential U = 5x2 + 10 ergs/gm. Find the
frequency.
Solution Given U = 5x2 + 10 ergs/g and m = 100 g
dU
F=- = -10 x
dx
d 2x
F=m = -10 x
dt 2
d 2x 10
or 2
=- x (i)
dt m
d 2x
Now = -w 2x (ii)
dt 2
Comparing Eqs. (i) and (ii),
10 10
w2 = or w =
m m
2p 10 1 10 1 10
= or n = =
T m 2p m 2 ¥ 3.14 100
n = 0.05 Hz
E xamplE 10 A lift is ascending at acceleration of 3 m/s2. What is the period of oscillation of simple pendulum
of length one meter suspended in the lift?
Solution Given f = 3 m/s2 and l = 1.0 m
The lift is ascending with an acceleration 3 m/s2 and acceleration due to gravity g = 9.8 m/s2. Hence, total acceleration
is 9.8 + 3 = 12.8 m/s2
l 1.0
Time period T = 2p = 2p
g¢ 12.8
= 1.755 s
T = 1.76 s
E xamplE 11 A mass of 6 kg stretches a spring 0.3 m from its equilibrium position. The mass is removed and
another body of mass 1.0 kg is hanged from the spring. What would be the period of motion if the spring is
now stretched and released?
Solution
F mg 6 ¥ 9.8
F = kx, k = = =
x x 0.3
k = 196 N/m
m 1.0
T = 2p = 2 ¥ 3.14
k 196
= 0.45 s
E xamplE 12 A spring when compressed by 10 cm develops a restoring force of 10 N. A body of mass 4 kg is

attached to it. Calculate the compression of the spring due to the weight of the body and calculate the period
of oscillation.
F 10
Solution Here F = kx or k = = = 100 N/m
x 0.1
The force applied F ¢ = mg = 4 ¥ 9.8 = 39.2 N
F ¢ 39.2
x¢ = = = 0.392 m
k 100
m 4
Time period T = 2p = 2 ¥ 3.14
k 100
T = 1.26 s
E xamplE 13 The relaxation time for damped harmonic oscillator is 50 s. Determine the time in which the
amplitude and energy of oscillator falls to 1/e times of its initial value.
Solution The amplitude of dampled harmonic oscillator at time t is given by
A(t) = A0e–st
1
Relaxation time t=
2s
given t = 50 s
1 1 1
Now t = = = per s
2s 2 ¥ 50 100
A0 is the amplitude at t = 0 and at time t the amplitude will be A0/e. Hence
A0 1
= A0e - st fi = e –st fi -1 = - st
e e
1
or t = = 100 s
s
E xamplE 14 Considering quality factor of sonometer wire of frequency 260 Hz as 2000, calculate the time
in which the amplitude decreases to 1/e2 of its initial value.
Solution The quality factor is given by
Q = wt
Here Q = 2000 and w = 2pn = 2 ¥ 3.14 ¥ 260 rad/s
Q 2000
Relaxation time t = =
w 2 ¥ 260 ¥ 3.14
= 1.225 s
The formula for amplitude of damped oscillator at time t is
A(t) = A0e–st
A
Given A(t) = 20
e
A0 A0
\ =
e 2 est
2
or t= = 2t
s
= 2 ¥ 1.225 = 2.450 s
Q1. Which one of the following is not necessary for SHM?

(a) elasticity (b) inertia (c) restoring force (d) gravity
Q.2 The velocity of a particle making SHM is maximum at
(a) equilibrium position (b) position of 1/3rd of amplitude
(c) extreme position (d) mid way (between zero and maximum)
Q.3 For a particle making SHM the phase difference between displacement and velocity is
(a) 0 (b) p (c) 2p (d) p/2
Q.4 Total energy of the particle executing SHM is same at
(a) Equilibrium position (b) midway (between zero and maximum)
(c) extreme position (d) all of them
Q.5 The potential energy of a particle executing SHM is maximum at
(a) equilibrium position (b) extreme position
(c) mid way (between zero and maximum) (d) position of 1/3rd of amplitude
Q.6 The kinetic and potential energies of a particle executing SHM are same at displacement (from mean
position)
(a) A (b) A/2 (c) A / 2 (d) A / 3
Q.7 The periodic motion which is not oscillatory, is
(a) simple pendulum (b) compound pendulum
(c) acoustic harmonic oscillator (d) motion of earth around sun
Q.8 The amplitude of a simple harmonic oscillator is doubled. The new time period would be
(a) same (b) double (c) half (d) none of these
Q.9 The unit of spring constant in SI system of units is
(a) Nm2 (b) Nm–1 (c) Nm–2 (d) Nm
Q.10 In which of the following oscillations the amplitude varies with time
(a) damped oscillator (b) forced oscillator
(c) undamped oscillator (d) none of these
Q.1 What is periodic motion?

Q.2 What do you mean by simple harmonic motion?
Q.3 Write the expression for the average energy of simple harmonic oscillator.
Q.4 State and explain the equation for simple harmonic motion.
Q.5 What is the time period for simple pendulum executing SHM and how it is modified for an harmonic
motion?
Q.6 What are the dimensions of forced constant of vibrating spring?
Q.7 Name the periodic motion which is not oscillatory.
Q.8 Are all the periodic motions simple harmonic? Is the reverse true? Explain.
Q.9 The amplitude of a simple harmonic oscillator is doubled. How does this effect the time period, total
energy and maximum velocity of the oscillator?
Q.10 What is damping? On what factors the damping depends?
Q.11 What is the effect of damping on the natural frequency of an oscillator?
Q.13 What do you understand by ‘quality factor’?
P ractice P roblems
general Questions
Q.1 Differentiate simple harmonic motion and oscillatory motion. Define simple harmonic motion.
Q.2 Derive a general differential equation of motion for a simple harmonic oscillator and obtain its solution.
Q.3 Simple harmonic motion is called sinusoidal or co-sinusoidal. Justify.
Q.4 Discuss the characteristics of simple harmonic oscillations. What is quality factor and how do you
define it?
Q.5 Derive an expression for the total energy of a harmonic oscillator and show that it is constant and
proportional to the square of the amplitude.
Q.6 When the displacement is one half of the maximum amplitude, what fraction of the total energy is
kinetic and what fraction is potential in simple harmonic motion?
Q.7 Derive a relation between restoring force of a spring and potential energy.
Q.8 Show that for a particle executing SHM the average values of kinetic and potential energies are the
same and each is equal to half of the total energy.
Q.9 Define damped harmonic oscillations. Solve the differential equation and discuss the case of oscillatory
motion. What is the quality factor and how do you define it?
Q.10 Discuss the theory of forced harmonic oscillations. How does sharpness of resonance depend on damping?
Q.11 What are damped vibrations? Establish a differential equation, of motion for a damped harmonic
oscillator and obtain an expression of displacement. Discuss the case of heavy damping.
Q.12 Define damped harmonic oscillations. Write the differential equation for a damped harmonic oscillator.
Solve the differential equation and discuss special cases of oscillatory motion.
Q.13 Write down the equation of damped simple harmonic oscillator. Find the expression for displacement
and discuss when we get oscillatory damped simple harmonic motion.
Q.14 Discuss the methods (logarithmic decrement, relaxation time and quality factor) for quantitative
measurement of damping effect in a damped simple harmonic oscillator.
Q.15 Explain free vibrations, damped vibrations, forced vibrations and resonance, giving one example of
each.
Q.16 Write note on
(i) harmonic oscillator (ii) Forced oscillations
Simple Harmonic Motion
and Sound Waves
8
Learning Objectives
LO 1 Understand superposition of two derivation of sound speed and intensity
simple Harmonic Motions (SHMs) of sound, and the level of sound
LO 2 Know about sound wave, its velocity, intensity
and sound displacement LO 5 Know about interference of sound
LO 3 Learn basics of standing waves, node, and waves in time (beats), and relation
anti-node along with detailed description between displacement and pressure
of their formation in air columns amplitude
LO 4 Understand Doppler effect, supersonic LO 6 Learn Lissajous figures and endoscopy
waves and shock waves along with together with its kinds
Introduction
Any periodic or oscillatory motion where the restoring force is proportional to the displacement and acts
opposite to the displacement is called a Simple Harmonic Motion (SHM). A simple example is the weight
attached to one end of a spring, the other end being tied to a rigid support such as a wall. If the mass is
displaced from the mean (equilibrium) position, the spring exerts a restoring force, according to the Hooke’s
law, F = –kx, where k is the spring constant of the spring and x is the displacement from the mean position.
According to the superposition principle, the net resultant at a given time and space, when two or more
SHMs combine, is the sum of the resultants by each of the SHMs. In physics, a standing wave, also known as a
stationary wave, is a wave in a medium in which each point on the axis of the wave has a constant amplitude.
A vibration that propagates as a typically audible mechanical wave of pressure and displacement through
a transmission medium such as air or water is called sound. The change in frequency or the wavelength of a
sound wave when the observer moves relative to the source, is termed as the Doppler effect. A common
example is the change in the pitch of the sound when the source moves towards or away from the observer
and vice-versa. A beat is an interference pattern between two sounds of slightly different frequencies. The
beat frequency is equal to the difference between the interfering frequencies.
8.1 SuperpOSitiOn Of twO SHMS LO1

For this, we consider the superposition of two Simple Harmonic Motions (SHMs), which produce a
displacement of the particle along the same line. We assume the case when both the motions have the same
frequency. The displacement of the particle produced by each simple harmonic motion is given by
y1 = a1 cos wt (i)
y2 = a2 cos (wt + d) (ii)
Here, d is the phase difference between y1 and y2.
The resulting displacement of the particle can be calculated by the linear combination
y = y1 + y2 = a1 cos wt + a2 cos (wt + d) (iii)
If d = 0, i.e., the two motions are in phase, then the resultant motion is
y = a1 cos wt + a2 cos wt = (a1 + a2) cos wt (iv)
This relation shows that the resultant motion is also an SHM with the same angular frequency w. However,
the motion has an amplitude equal to the sum of the amplitudes of the two motions, i.e., a = a1 + a2.
8.2 SOund wave and itS veLOcity LO1

Let us look at the mechanism of sound-wave generation in air. We exert a pressure on the air molecules
while we speak because of which compression of these molecules takes place (i.e., they come closer to each
other). Then these molecules try to come back to their original position but cross the equilibrium position,
due to which rarefaction is produced. In the next instance, these rarefactions are converted into compressions.
Hence, the longitudinal wave in air is generated where the molecules execute SHMs. This is the way sound
waves propagate in air. In general, we can say that the sound wave is the propagation of pressure variations
in elastic media. Both the longitudinal sound waves and transverse sound waves, however, occur in solid
elastic media. In a three-dimensional homogeneous (uniform) medium, the sound from an ideal point source
propagates in the form of spherical waves.
The sound velocity depends on the properties of the medium, but for the large amplitudes, this also depends
Ê CP ˆ
on the amplitude. The sound velocity in gases (say, Csg) depends on the adiabatic coefficient g Á ∫ ,
Ë CV ˜¯
density r and the pressure p or the temperature T of the gas. This is given by
gp
Csg = = g RsT (i)
r
Here, Rs is the specific gas constant in J/K kg. In most gases, the sound velocity is in the range of
200–1300 m/s.
The sound velocity in liquids, (say, Csl) depends on the compression modulus K (measured in N/m2) and the
density r of the liquid. Specifically, it is given by
K
Csl = (ii)
r
Simple Harmonic Motion and Sound Waves 261
In most liquids, the sound velocity is in the range of 1100–2000 m/s, and in water it is 1480 m/s at 20°C
temperature.
The sound velocity in solids (say, Css) depends on the elasticity modulus E (measured in N/m2) and the
density r of the solid. Specifically, it is given by
E
Css = (iii)
r
As mentioned earlier, sound waves in solids may be longitudinal waves or transverse waves. The sound
velocity in most solids is in the range of 1200–6000 m/s, and in iron, it is 5000 m/s.
8.3 SOund diSpLaceMent LO1

So far, we have understood that the particles of the medium execute SHM in order to generate the sound
wave, which means they vibrate about their equilibrium position. Sound displacement is defined as the
displacement of the vibrating particles of the medium from their rest positions. This can be represented by
y(x, t) = y0 sin (wt – kx) (i)
where w is the angular frequency of oscillations (w = 2pf) and k is the wave number (k = 2p/l, l is the
wavelength). If we differentiate y w.r.t. time t, we get the velocity of vibrating particles, which is also known
as sound-particle velocity.
This is obtained as
dy
v= = y0w cos(w t - kx )
dt
or v( x, t ) = y0 w cos(w t - kx ) (ii)
This equation shows that there is a phase difference of p/2 between y and v. Moreover, the velocity v is
different for different sound particles.
8.4 Standing waveS LO1

Standing waves are produced by superposition of two waves of equal frequency, amplitude and phase, while

they propagate in opposite directions. It means the wave vectors k of these waves should be anti-parallel. If
we consider the superposing waves to propagate in x and –x directions, these can be represented by
y1(x, t) = a cos (kx – wt) (i)
y2(x, t) = a cos (–kx – wt) (ii)
The resultant wave achieved after the superposition is obtained as
y(x, t) = y1 + y2 = –2a cos wt cos kx (iii)
The above equation represents the standing wave, whose displacement y(x, t) is shown in Figure 8.1 as a
function of x for different instants of time t. Here, it is seen that the minima and maxima of the standing wave
are fixed in space, unlike the other propagating or travelling waves. The notation for a space-fixed minimum
of a standing wave is termed node, while that for a space-fixed maximum is known as anti-node.
Figure 8.1
8.5 Standing waveS in air cOLuMnS LO1

Standing waves can be set up in a tube of air, for example an organ pipe, by the superposition of longitudinal
sound waves travelling in opposite directions. Here, the incident wave can be reflected from one end of the
pipe. The phase relationship between the incident and reflected waves depends on whether that end is closed
or open.
In a pipe closed at one end, a displacement node will always be found at the closed end because the wall
at this end does not allow longitudinal motion of the air molecules. It also means that the reflected sound
wave is 180° out of phase with the incident wave at a closed end of the pipe. The closed end of such an air
column corresponds to a pressure anti-node (a point of maximum pressure variation), since the pressure wave
is 90° out of phase with the displacement wave. On the other hand, a pressure node or displacement anti-
node appears approximately at the open end of the air column. At the open end, pressure variation does not
occur as this end of the air column is open to the atmosphere, where the pressure must remain constant at the
atmospheric pressure.
Standing waves arise in the air column (length L) in a pipe with its one end closed, if the wavelength ln
satisfies the following condition, which is also known as resonance condition.
4L
ln = (i)
2n + 1
The frequencies corresponding to ln are referred to as natural frequencies fn. For the sound waves (velocity
Cs), these are given as
Cs Cs (2 n + 1)
fn = = (ii)
ln 4L
Fundamental vibration corresponds to the standing wave with n = 0 (n is a non-negative integer). It means
the fundamental wavelength l0 and the fundamental frequency f0 are obtained as
Cs
l0 = 4 L, f0 = (iii)
4L
Cs st
l 0 = 4 L , f0 = 1 harmonic
4L
4L 3C
l1 = , f1 = s = 3 f0 3rd harmonic
3 4L
4L 5C
l2 = , f2 = s = 5 f0 5th harmonic
5 4L
Figure 8.2(a)
A standing wave with a node number n different from zero (n > 0) is called the harmonic. It is clear that the
natural frequencies or the places of nodes and anti-nodes depend on the length of the pipe, L. The standing
waves in a pipe closed at one end are shown in Figure 8.2(a). Here, it can be seen that the natural frequencies
of oscillation form a harmonic series that includes only odd integral multiplies of the fundamental frequencies.
Cs st
l 0 = 2 L , f0 = 1 harmonic
2L
Cs
l1 = L, f1 = = 2 f0 2 nd harmonic
L
2L 3Cs
l2 = , f2 = = 3 f0 3rd harmonic
3 2L
Figure 8.2(b)
In an open pipe (two ends free), the resonance condition reads

2L C (n + 1)
ln = or fn = s (iv)
n +1 2L
It means the fundamental variation has wavelength l0 and frequency f0, given below.
Cs
l 0 = 2 L , f0 = (v)
2L
The first three normal modes of oscillation of a pipe open at both ends are shown in Figure 8.2(b).
8.6 dOppLer effect LO1

Doppler effect is the phenomenon that gives rise to a change in the observed wavelength or frequency of
sound or radiation, which results from the movement of its source relative to the observer. A source moving
towards the observer appears to decrease the wavelength or to increase the frequency, whereas a source
moving away from the observer appears to increase the wavelength or to decrease the frequency. The change
in wavelength is known as the Doppler shift.
The Doppler effect is widely used to measure velocities by reflection of a wave transmitted from the moving
object. This effect is also used to measure the velocities of distant galaxies based on the observed red shift.
Not only this, ultrasound for blood movement in arteries makes use of the Doppler effect.
8.6.1 Moving Source
Consider a source of sound wave to be at rest (Figure 8.3). The wave
crests corresponding to the emitted sound wave can be represented by
circles whose centre is at the position of the source. If the frequency
of sound is f0, then these crests are generated at the frequency f0 only,
and the separation between successive crests will be the wavelength l0
of the sound. In view of Cs as the sound speed, f0 and l0 are related to
each other as follows:
Cs Figure 8.3
l0 = = C st 0 (i)
f0
Here, t0 is the time interval at which the observer receives these successive wave crests.
Now, consider the source S to move towards the observer at the speed vs, which is less than Cs, i.e., vs <<
Cs. Then in time t0, the source will cover a distance vst0 towards the observer. At the time t0, the previously
emitted crest will itself have moved towards the observer by a distance l0. Hence, the actual distance between
the successive crests emitted towards the observer will be
l¢ = l0 – vs t0 (ii)
Corresponding to the wavelength l¢, the observer will observe the frequency of sound as
Cs Cs C s /l 0
f¢= = =
l ¢ l0 - vst 0 (1 - vst 0 /l0 )
f0
or f¢= (iii)
(1 - vs /Cs )
v
For the common case of s << 1 , Eq. (iii) can be
approximated as Cs
Ê v ˆ
f ¢ = f0 Á 1 + s ˜ (iv)
Ë C s¯
The above observation is picturized in Figure 8.4.

Figure 8.4
Based on the above argument, the frequency observed by an observer in the case of the source moving away
from this would be
f0
f¢= (v)
(1 + vs /Cs )
8.6.2 Stationary Source and Moving Observer

We can also discuss the case when the source of
the sound wave remains at rest but the observer
moves towards it at the speed vobs. It is better
to recall here that the sound-wave crests are l0
apart in the air and moving at the speed Cs. Let
us assume that the observer meets the successive
crests in time t (different from t0). In this time,
the observer moves a distance vobst and the crests
also move towards the observer by a distance Cst Figure 8.5
(Figure 8.5).
From Figure 8.5, it is clear that
l0
t= (vi)
Cs + vobs
This gives the frequency of sound observed by the observer as

1 Cs + vobs
f = =
t l0
Cs Ê vobs ˆ
or f = ÁË 1 + C ˜¯
l0 s
Ê v ˆ
= f ¢ Á1 + obs ˜ (vii)
Ë Cs ¯
Finally, we can also discuss the case when both the source and observer move towards each other, keeping in
view Eqs (vii) and (iii), the observed frequency can be obtained as
(1 + vobs /Cs )
f = f0
(1 - vs /Cs )
(Cs + vobs )
or f = f0 (viii)
(Cs - vs )
It means the frequency observed in the case of a moving source and observer towards each other is larger than
the frequency observed when either source or observer moves alone and approach the other.
8.7 SuperSOnic and SHOck waveS LO1

In the case when the speed of a source is equal to the speed of sound, the wavefronts produced by the source
cannot escape it. Hence, the waves will pile up and form a large amplitude sound barrier. This barrier makes
persistent flight at this speed not only difficult but also risky. We will explain this phenomenon below.
Let us represent the speed of source, for example, an aeroplane, as v relative to the air. If the speed v is less
than the speed of sound Cs, the waves in front of the aeroplane will be crowded together with a wavelength
given by
Cs - v (i)
l=
f
This equation shows that the wavelength approaches zero as the speed of the aeroplane approaches the speed
of sound. Then the wave crests will pile up on each other, as shown in Figure 8.6(a) According to Newton’s
third law, now the aeroplane must exert a large force to compress the air in front of it and the air exerts an
equal and opposite large force on the aeroplane. It means there is a large enhancement in aerodynamic drag
or the air resistance, as the aeroplane approaches the speed of sound. This phenomenon is known as the sound
barrier.
Figure 8.6
On the other hand, when the speed v of the source is greater than the speed Cs of the sound, the source of the
sound, the aeroplane, is called supersonic. The aeroplane during its motion produces sound by displacing the
surrounding air, and a series of wave crests is emitted from the front side (nose) of the aeroplane. Each of the
wave crests spreads out in a circle centred at the position of the aeroplane. After a time t, the crest emitted
from the initial position of the aeroplane (say, S1) spreads to a circle of radius Cst, whereas the aeroplane
moves a greater distance vt (say, to the position S2) in view of v > Cs. Under this situation, it can be seen
that the circular crests interfere constructively at the points along the line which makes an angle q with the
direction of the velocity of the aeroplane. This leads to a very large amplitude wave crest along this line. This
large amplitude crest is known as a shock wave. It is also evident that a shock wave forms a cone around the
direction of motion of the source.
We can calculate the angle q from Figure 8.6(b), where the right-angled triangle (DS1NS2) shows that
Cs t Cs
sin q = = (ii)
vt v
In view of v > Cs, the ratio v/Cs is called the Mach number, which is greater than unity for all supersonic
speeds. Until the source such as a supersonic jet aeroplane or a rifle bullet moves with constant velocity, the
angle q remains a constant, and the shock-wave cone moves along with the source only. Here, it will be worth
mentioning that the sonic boom we hear after a supersonic aeroplane has passed by is the arrival of this shock
wave only. Unlike ordinary sound waves, the speed of a shock wave depends on its amplitude. The shock
wave speed is always greater than the sound speed in the fluid, and it decreases as the amplitude of the wave
reduces. This also means that the shock wave will die and reduce to an ordinary sound wave when the speed
of the shock wave equals the normal speed of the sound.
8.8 derivatiOn Of SOund Speed LO1

Let us take the example of sound waves in air or neutral gas. We represent the mass density by r, pressure

by p and the velocity of air molecules by v . In order to derive the sound speed, we can use the concept of
conservation of mass (continuity equation) and the equation of motion (Navier–Stokes equation). Considering
adiabatic compression, we can write the pressure in terms of density r, temperature T and the ratio g of
specific heats at constant pressure Cp and at constant volume Cv.
—p —r
=g (i)
p r
With this, the equation of motion reads

∂v —r
 
r = -—p = -g p (ii)
∂t r
The continuity equation is written as

∂r 
+ — ◊ (rv ) = 0 (iii)

∂t
Ê ∂ˆ
For the one-dimensional case Á — = x̂ ˜ , the linearized form of Eqs (ii) and (iii) can be found by taking
Ë ∂x ¯
r = r0 + r1 and v = v1. Here, the quantity r0 is the density in equilibrium and r1 is due to oscillations or
generation of wave. Similarly, v1 is the velocity of air molecules during oscillations. Hence, we get
∂v1 g p ∂r
r0 =- 0 1 (iv)
∂t r0 ∂ x
∂r1 ∂ (v)
+ ( r0 v1 ) = 0
∂t ∂x
For the wave of frequency w and wave number k, the oscillating quantities can be taken to have dependence as
r1 = r1 ei(kx – wt) (vi)
v1 = v1 ei(kx – wt) (vii)
∂ ∂
These equations show that can be replaced with ik and with –iw. So Eqs (iv) and (v) read
∂x ∂t
g p0
-iwr0 v1 = - ik r1 (viii)
r0
–iw r1 + r0 ikv1 = 0 (ix)
We solve Eqs (viii) and (ix) for v1 and obtain

Ê 2 2 g p0 ˆ (x)
ÁË w - k r ˜¯ v1 = 0
0
Now, for the occurrence of wave, v1 π 0.

g p0
Hence, w 2 - k 2 =0
r0
w g p0
or = (xi)
k r0
w
The ratio is the phase velocity of the wave, which we call sound speed (say, Cs). In terms of temperature
k
T of the gas or air, this can be written as
w g KTn0
= Cs = (n0 is the number density)
k n0 M
g KT
fi Cs = (xii)
M
Here, M is the mass of each air molecule.
8.9 intenSity Of SOund and itS derivatiOn LO1

Intensity of sound at a point in a progressive wave is defined as the sound energy per unit area per unit time
perpendicular to the direction of propagation of the wave. It is measured in W/m2 in the SI system of units.
We consider a plane progressive simple harmonic wave travelling along the positive x-direction with velocity
v(= w/k). The displacement y at a time t can be represented as
y = a sin (wt – kx) (i)
From this, we find the velocity of particle by differentiating it w.r.t. time
dy (ii)
= w a cos (w t - kx )
dt
In order to calculate the energy or the intensity of sound, we consider a medium of density r. Taking unit area
of medium having thickness dx perpendicular to the direction of propagation of wave, we find the kinetic
energy as
2
1 Ê dy ˆ
dK = r dx Á ˜ (iii)
2 Ë dt ¯
dy
Putting the value of from Eq. (ii) and using w = 2pf, where f is the linear frequency of the wave, Eq. (iii)
dt
reads
1 2
dK = r a (2p f )2 cos2 (w t - kx ) dx
2
= 2p2a2rf 2 cos2 (wt – kx) dx (iv)
This shall give the total energy of the wave as
dE = dKmax (when potential energy is zero)
= 2p2 a2rf 2 dx (v)
dx can be written in terms of the velocity v as dx = vdt.
Hence,
dE = 2p2a2f 2r v dt (vi)
The integration gives
E = 2p2a2f 2r vt (vii)
2 2 2
The energy flow per unit time is obtained from Eq. (vii) as 2p a f rv, which is nothing but the intensity of
the sound wave. Hence,
I = 2p2a2f 2r v (viii)
The other forms of the formula of sound intensity I are
I = 2p2rf 2 v A2max (in terms of displacement)

2
r v Dvmax
I=
2
(in terms of velocity, where Dvmax = 2p f Amax )
2
Dpmax
I=
2rv
(in terms of pressure, where Dpmax = 2pr fv Amax )
8.10 SOund-IntenSIty LeveL LO1

The level of sound intensity (say IL) can be defined in terms of decibels (dB) and neper (Np), as follows:
Ê I ˆ
I L = 10 log10 Á ˜ in dB (i)
Ë I0 ¯
Ê I ˆ
or I L = log10 Á ˜ in B (Bel)
Ë I0 ¯
1 Ê I ˆ
IL = log e Á ˜ in Np (ii)
2 Ë I0 ¯
Here, I is the sound intensity, I0 is the reference intensity and B is the unit bel (1 B = 10 dB).
Since the sound intensity I is directly proportional to square of the pressure p, we have
I p2
= 2 (p0 is the reference pressure)
I 0 p0
With the help of this, IL can be defined in terms of p as follows:

Ê p2 ˆ
I L = 10 log10 Á 2 ˜
Ë p0 ¯
Ê pˆ
= 20 log10 Á ˜ (iii)
Ë p0 ¯
8.11 interference Of SOund waveS in tiMe: BeatS LO1

A beat is an interference between two sound waves having slightly different frequencies. So beating is
obtained when two sound waves of different frequencies superimpose with each other. Under this situation,
alternating constructive and destructive interference cause the sound to be alternatively loud and soft. This
phenomenon is called beating or producing beats.
The beat frequency is given by the difference of frequencies of respective waves, i.e.,
1
fbeat = f1 - f2 = (i)
Tbeat
Let us consider two sound waves having different frequencies (w1 and w2), which superimpose. Then, the
displacement perceived will be given by
h(t) = sin w1t + sin w2t
Ê w + w2 ˆ Ê w1 - w 2 ˆ
= 2 sin Á 1
Ë ˜¯ t ◊ cos ÁË ˜t (ii)
2 2 ¯
w1 + w 2 w - w2
If we represent = w 0 as the average frequency and 1 = Dw as the difference in frequencies,
2 2
then
h(t) = 2 sin (w0t). cos (Dwt) (iii)
Clearly, the time period
2p 2p 1
Tbeat = = =
Dw w1 - w 2 f1 - f 2 (iv)
2
Hence, the beat frequency
1
fbeat = = f1 - f2 (v)
Tbeat
If f2 > f1, then fbeat = f2 – f1. This is the reason the modulus of the difference is taken in Eq. (i).
8.12 reLatiOn Between diSpLaceMent and preSSure aMpLitude LO1

Consider a plane monochromatic wave of sound travelling with a phase velocity v0 along the length of the
tube filled with a gas with density r0. We will find the relationship between amplitude of the molecular
displacement and that of pressure oscillations.
Proof
Figure 8.7
Consider a long tube of cross-sectional area A, in which a movable piston is fitted easily in the left side.
The tube is open at another end with ambient surrounding. Firstly, the piston is at rest and then we apply a
force F on the piston so that it moves towards the right side. This way we produce compression so that the
wavefront moves dx distance towards right during the small interval dt. If the wavefront moves with velocity
v0, then
dx = v0dt (i)
Due to the moment or displacement of the piston, the gas molecules also move with lower velocity vm. Since
the distance moved by the piston is very small, we can consider that the molecules under that volume have
the same speed as that of the wavefront. So the total mass moved by the gas is obtained as
dm = r0 A dx (ii)
Here, r0 is the mass density of the gas.
Linear momentum is given by
dP = dm · vm
= r0 A dx vm (iii)
Putting the value of dx from Eq. (i), we get
dP = r0 A vm v0 dt (iv)
From Eq. (iv), we can obtain the force as follows:

dp
F=
dt
1 dp
and the pressure P =
A dt
= r0 v0 vm (v)
If the maximum longitudinal speed is sw, where s is the displacement of molecules of the gas and w is the
angular frequency of the sound wave, then maximum pressure is given by
Pmax = r0 v0 sw (vi)
This is the required relation between the displacement s and the pressure amplitude P or Pmax. Clearly, a
larger pressure is exerted by a larger velocity of the piston (v0) or the larger displacement of the molecules
(s) of the sound.
8.13 LiSSajOuS figureS LO1

Jules Antoine Lissajous in 1857 demonstrated that when a particle is acted upon simultaneously by two
simple harmonic motions at right angles to each other, the resultant path traced out by the particle is a curve.
The curves so obtained are called Lissajous figures. So, Lissajous figures are the patterns formed when two
vibrations transverse to each other are superimposed. In other words, a Lissajous figure is a parametric plot
of a harmonic system.
Let us consider that each vibration is a simple harmonic motion represented by a sinusoidal wave. Their
amplitudes are a and b and frequencies are w1 and w2; f1 and f2 are their phases and t is the time. We plot an
input x(t) and output y(t) of a linear time-independent system.
x = a sin (w1t + f1), y = b sin (w2t + f2) (i)
After time t, the point P of the coordinates (x, y) will trace a curve. The equation of the curve can be found by
eliminating t in Eq. (i). For example, when both the frequencies are same and there is no phase difference, i.e.,
w1 = w2 = w and f1 = f2 = 0°, then we have
x = a sin wt, y = b sin wt
x y b
fi = and thus y = x (ii)
a b a
This represents an equation of a straight line. It means when two sine waves
are of equal frequency and are in phase, then we get a diagonal line to the
right, as shown in Figure 8.8.
If the frequencies are same and the two vibrations are out of phase, i.e., w1
= w2 = w and f1 – f2 = p, then Eq. (i) gives
b Figure 8.8
y=- x (iii)
a
This is also an equation of a straight line. It means when two sine waves are of equal frequencies and 180°
out of phase, we get a diagonal line to the left, as shown in Figure 8.9.
Now, we consider the case of p/2 phase difference, i.e., w1 = w2 = w and

f1 – f2 = p/2. Under this case, Eq. (i) gives
x = a sin wt, y = b sin (wt + p/2) = b cos wt (iv)
Equation (iv) yields
x2 y2
= sin 2 w t and = cos2 w t
a2 b2
which on adding gives
x2 y2 (v)
+ =1
a2 b2 Figure 8.9
The above equation represents an ellipse as shown in Figure 8.10.
If a = b, then, Eq. (v) represents a circle x2 + y2 = a2, as shown in Figure 8.11.
It means the circular path (Lissajous figure) is obtained when two sine waves are 90° out of phase
superimposing each other.
Finally, it can be seen that if the phase difference varies continuously, y(t)
then the ellipse will slowly change its orientation and shape.
This is shown in Figure 8.12.
The curve becomes more complex when we consider the case of unequal b
frequencies. If we consider w2 = 2w1 and f1 – f2 = p/2, then Eq. (i) reads
a x(t)
x = a sin wt, y = b sin (2wt + p/2)
= b cos 2 wt
= b[1 – 2 sin2wt] (vi)
From Eq. (vi), we get Figure 8.12
È Ê xˆ
2˘
y = b Í1 - 2 Á ˜ ˙
ÍÎ Ë a¯ ˙˚
This equation represents a parabola. Following curves can also be obtained based on different frequencies
and phase differences.
Figure 8.13
If we know both the angle of the major axis of a Lissajous curve and the direction of the curve’s rotation,
then we can determine the quadrant of phase shift d ◊ (∫ f1 – f2).
Taking f1 – f2 = d for the case of w1 = w2 = w, we can summarize the above results as below:
d = 0° Line with positive slope
0° > d > –90° Curve in counterclockwise direction with positive slope
d = –90° Counterclockwise circle
–90° > d > – 180° Counterclockwise curve with negative slope
d = –180° Line with negative slope
–180° > d > –270° Clockwise curve with negative slope
d = –270° Clockwise circle
–270° > d > –360° Clockwise curve with positive slope
8.14 endOScOpy LO1

Endoscopy means typically looking inside the human body for medical reasons. This is done by an endoscope.
It uses an optical-fibre technology that dates back to the 1880s when Dr. Roth and Prof. Reuss of Vienna used
a bent glass rod to illuminate the body cavity. Based on endoscopy, many medicines have been revolutionized,
since these have the capacity of reaching inside the human body in the least harmful way. All the endoscopes
make use of fibre-optic technology to illuminate their images.
Fibre-optic endoscopes are flexible and highly maneuverable instruments that allow access to channels in the
body, which older semi-rigid instruments cannot access at all or can access only at great discomfort to the
patient. Composed of multiple hairlike glass rods bundled together, these instruments can be more easily bent
and twisted. Finally, the intense light enables the endoscopist to see around.
8.14.1 types of endoscopes
There are two major types of endoscopes, namely, rigid endoscopes and flexible endoscopes. The details of
these are given below.
Rigid Endoscopes
Also known as a laparoscopes, these are basically medical periscopes. A laparoscope is a long fibre-optic
cable system that allows viewing of the affected area by inserting the cable from a more distant but more
accessible location. There are a number of advantages to the patients having laparoscopic surgery over the
open procedure. For example, pain and haemorrhaging are reduced due to smaller incisions. On top of it, the
recovery times are also shorter.
Flexible Endoscopes
These are based on either fibre optics or LCD. The use of flexible endoscopes is common in both medical and
surgical specialities. These endoscopes provide the unique ability to reach cavities and viscera, which are not
visible to the naked eye. These allow for minimally invasive investigation of symptoms, diagnosis pathology
and application of directed therapies. Advances in the imaging systems, newer endoscopes with ‘self-drive’
capabilities and enhancement of targeted therapeutics are future applications of flexible endoscopy.
s UmmarY
✦ The principle of superposition states that when two waves of the same kind meet at a point in space,
the resultant displacement at that point is the vector sum of the displacements that the two waves would
separately produce at that point. Superposing of two or more coherent waves to produce regions of
maxima and minima in space results in the production of interference pattern. Constructive interference
occurs when two or more waves arrive at the screen in phase (phase difference 0 or 360 degree) with
each other, so that the resultant wave amplitude is the sum of the amplitude of the individual waves.
Destructive interference occurs when the two or more waves arrive p out of phase with each other.
✦ Simple harmonic motion is any motion where a restoring force is applied that is proportional to the
displacement and in the opposite direction of that displacement. In other words this also means that the
acceleration is proportional to displacement but they are in opposite directions.
✦ Standing waves or stationary waves are produced when two different waves of same kind having same
frequency, amplitude, and phase, superimpose when travelling in opposite direction. It results in the
formation of regions of maxima and minima. The region of minimum or zero displacement is called a
node, while the region of maximum displacement is called an antinode.
✦ Standing waves arise in a number of situations, for example an air column. The one end of the pipe has to
be closed for waves to occur. The wavelength of the wave must satisfy a condition called the resonance
condition, which is a function of ‘n’. The frequency corresponding to different wavelengths are called
natural frequency. Standing wave corresponding to n = 0 corresponds to Fundamental vibrations.
✦ Doppler Effect is the phenomena which results in the change in the frequency of the sound waves
whenever there is relative motion between the source and the observer. The corresponding change in
the frequency or wavelength is called Doppler Shift. Doppler shift is used to measure the velocities of
distant galaxies based on their recessional red shift. This has led to the observation that the universe is
expanding.
✦ In Doppler Effect, it is not the motion of the individual source or the observer that matters, but the
relative motion between the source and the observer.
✦ When the speed of a source equals the speed of sound (v = c) the wave fronts cannot escape the source.
The resulting pile of waves forms a large amplitude “sound barrier” that makes sustained flight at this
speed difficult and risky. When the speed of a source exceeds the speed of sound (v > c) the wave fronts
lag behind the source in a cone-shaped region with the source at the vertex. The edge of the cone forms
a supersonic wave front with an unusually large amplitude called a “shock wave”. When a shock wave
reaches an observer a “sonic boom” is heard.
✦ A beat is an interference pattern between two sounds of slightly different frequencies. The beat
frequency is equal to the difference between the interfering frequencies.
✦ Lissajous figure, also called Bowditch Curve, are the patterns produced by the intersection of two
sinusoidal curves or waves which are at right angles to each other. First studied by the American
mathematician Nathaniel Bowditch in 1815, the curves were investigated independently by the French
mathematician Jules–Antoine Lissajous in 1857–58. So basically these are the patterns which are
formed when two sinusoidal waves interfere at right angle to each other.
✦ Lissajous figures can also be used as an experimental setup in laboratory to find the frequency of an
unknown source when frequency of one of the sources is known.
✦ Endoscopy means looking inside the human body for medical reasons which is done by endoscope.
Fibre optic endoscopes are flexible and highly maneuverable instruments which allow access to
channels in the body, which older semi rigid instruments cannot access at all or can access only at great
discomfort to the patient.
✦ Two types of endoscopes are generally used. These are named as rigid endoscope and flexible
endoscope. The rigid endoscope are also known as laparoscopes.
s olved e XamPles
E xamplE 1 A source of sound is travelling east at 10 m/s toward you. You are travelling at 2 m/s east. It is
20°C. When the source is not moving, it emits a sound of 3000 Hz frequency. What frequency do you hear?
Sound in air at 20°C travels at 343 m/s.
Solution Given us = 10 m/s, uobs = –2 m/s
f0 = 3000 Hz, v = 343 m/s
Ê 1 + uobs /v ˆ
f ¢ = f0 Á
Ë 1 - us /v ˜¯
Ê 1 - 2/343 ˆ
= 3000 Á = 3072 Hz
Ë 1 - 10 / 343 ¯˜
E xamplE 2 Suppose a train is approaching you while you are standing on the platform at the station. As the
train approaches the station, it slows down but the engineer is sounding the hooter at a constant frequency
of 400 Hz. Describe the pitch of the hooter and the changes in pitch of the hooter that you hear as the train
approaches you. Take the speed of sound in air as 340 m/s.
Solution The frequency of the sound gradually increases as the train moves towards you. The pitch increases and you
shall hear a higher pitched sound.
E xamplE 3 Passengers on a train hear its whistle at a frequency of 750 Hz. Ram is standing next to the train
tracks. If the train moves directly towards him at a speed of 30 m/s. what frequency does he hear? Take the
speed of sound in air as 340 m/s.
Solution fL = ((v + vL)/(v + vS)) ¥ fS = (340 + 0)/(340 – 30)) ¥ (750) = 822.58 Hz
E xamplE 4 A small aircraft is taxiing directly away from you down a runway. The noise of its engine, as the
pilot hears, has a frequency 1.20 times the frequency that you hear. What will be the speed of the plane? Take
the speed of sound in air as 340 m/s.
Solution The velocity of listener (you) is 0 and the source is moving away from you at an unknown velocity. This
velocity must be positive. We also know that
fS = 1.20 fL
fL = (v + vL)/(v + vS) ¥ fS = (340 + 0)/(340 + vS) ¥ (1.20 fL)
340 + vS = (340)(1.20)
This gives the velocity of the source vS = 68 m/s
E xamplE 5 Suresh is in his car moving at the speed of 0.50c towards Sheetal who is sitting in her stationary
car. It is getting dark and Sheetal does not have her headlights on, so Suresh flashes his brights at Sheetal. If
the frequency of the light which Suresh emits from his headlights is 4400 Hz, at what frequency does Sheetal
hear the sound?
Solution Since the source is moving at speed vS, the appropriate formula is f1 = f/(1 ± vS/v). Here f1 is the changed
frequency and f is the initial frequency. The changed frequency will be faster because Suresh’s car is coming towards
Sheetal. So the bottom of the fraction should be less than one. It means we should use the minus sign rather than the plus
sign. Therefore, the formula yields
f1 = (4400)/(1 – (.50c)/(c)) = 8800 Hz
E xamplE 6 Deepika is walking down the streets of downtown New Delhi and comes to an intersection
where the walk signal is blinking indicating to stop walking. Deepika thinks she is smarter than the signal
and tries to make a last minute run to the other side of the street. She realises that she is not going to make
it when a speeding truck coming towards her at the speed of 25 m/s is honking at the frequency of 6040 Hz.
With what frequency is the wave reaching Deepika right before she gets struck by the truck, if the speed of
sound is 343 m/s?
Solution Formula to be used is
f1 = f/(1 – vS/v)
This gives
f1 = 6040/(1 – 25/343) = 6514.86 Hz
E xamplE 7 As a train pulls out of the station going 60 m/s, it blasts its horn. What would be the frequency
heard by the passengers in train if the passengers still at the station are hearing 380 Hz?
Solution Since the frequency of the sound heard by the passengers at the station must be lower, the formula should be
f1 = f/(1 + vS/v)
This gives
380 = f/(1 + 60/343) or f = 446 Hz
E xamplE 8 Mukesh is on a motorcycle speeding down the highway at 45 m/s until he sees a traffic jam
ahead. The honking made by the stopped cars is 780 Hz, what frequency does Mukesh hear the sound at?
Solution Since the frequency Mukesh hears will be higher than what the light is actually emitting, the appropriate
formula is
f1 = f(1 + vS/v)
This gives
f1 = 780(1 + 45/343) = 882 Hz.
E xamplE 9 The speed of sound waves in air is found to be 340 m/s. Determine the fundamental frequency
(first harmonic) of an open-end air column that has a length of 67.5 cm.
First
harmonic
L = 0.675 m
Figure 8.14
Solution Given, v = 340 m/s, L = 0.675 m

For the first harmonic, wavelength is twice the length.
l1 = 2 L = 2 × 0.675 m = 1.35 m
v v 340 m/s
f1 = = = = 252 Hz
l1 2 L 1.35 m
A wave of 1000 Hz frequency travels in air of 1.2 kg m–3 density at 340 m/s. If the wave has 10
E xamplE 10
m Wm intensity, find the displacement and pressure amplitudes.
–2
1
Solution I= ( rv)( Aw )2 v Æ wave speed
2
2I
fi A =
r vw 2
2 ¥ 10 -6
=
1.2 ¥ 340 ¥ (2p ¥ 1000)2
= 11 nm
r0 = rv Aw
= 1.2 ¥ 340 ¥ 11 ¥ 10–3 ¥ 2p ¥ 1000
= 28 mPa
E xamplE 11 Assuming r = 1.29 kg/m3 for the density of air and v = 331 m/s for the speed of sound, find the
pressure amplitude corresponding to the threshold of hearing intensity of 10–12 W/m3.
1 2
Solution I= Pmax / r0 v
2
fi Pmax = 2 I r0 v
= 2 ¥ 10 -12 ¥ 1.29 ¥ 331

= 2.92 ¥ 10–5 N/m2
E xamplE 12 For ordinary conservation, the intensity level is given as 60 dB. What is the intensity of the
wave?
I
Solution IL = 10 log
I0
I
60 = 10 log
10 -12
log I + log 1012 = 6
log I = –6
\ I = 10–6 W/m2 = 1 mW/ m2
E xamplE 13 A small source of sound radiates energy uniformly at a rate of 4 W. Calculate the intensity level
at a point 25 cm from the source if there is no absorption.
Power 4
Solution I= 2
= = 5.093 ¥ 10–4 W/m2
4p r 4p ¥ 252
I
IL = 10 log
I0
5.093 ¥ 10 -4
= 10 log
10 -12
= 10 log (5.093 ¥ 108)
= 10[log 5.093 + 8]
= 87 dB
E xamplE 14 The maximum pressure variation that the ear can tolerate is about 29 N/m2. Find the
corresponding maximum displacement for a sound wave in air having a frequency of 2000 Hz. Assume the
density of air as 1.22 kg/m3 and the speed of sound as 331 m/s.
Pmax Pmax
Solution A= 2
=
k r0 v 2p r0 fv
k = 2p/l and v = fl
29
A= = 5.7 ¥ 10–6 m
2 ¥ 3.14 ¥ 1.22 ¥ 331 ¥ 2000
E xamplE 15 If two sound waves, one in air and the other in water, have equal pressure amplitude, what is
the ratio of intensities of waves? Assume that the density of air is 1.293 kg/m3, and the speeds of sound in air
and water are 330 and 1450 m/s respectively.
2
Pmax
Solution I=
2 r0 v
Pmax (air) = Pmax (water)
I Water r v
\ = A A
I Air rW vW
1.293 ¥ 330
= = 2.94 ¥ 10 -4
1000 ¥ 1450
E xamplE 16 The pressure in a progressive sound wave is given by the equation P = 2.4 sin p(x – 330 t),
where x is in metres, t in seconds and P in N/m2. Find (a) pressure amplitude, (b) frequency, (c) wavelength,
and (d) speed of wave.
Solution P = 2.4 sin p(x – 330 t)
Ê1 ˆ
= 2.4 sin 2p Á x - 165 t ˜
Ë2 ¯
Êx ˆ
P = Pmax sin 2p Á - ft ˜
Ël ¯
On comparing, we get
Pressure amplitude = 2.4 N/m2
Frequency = 165 Hz
Wavelength = 2.0 m
Speed of wave v = fl = 165 ¥ 2 = 330 m/s
Q.1 Damping of periodic motion always results in

(a) decreasing amplitude (b) energy gain in vibration
(c) resonance (d) energy loss to friction
(e) none of these
Q.2. If the frequency of a 0.75 m simple pendulum is 1.5 Hz, the angular frequency on a corresponding
reference circle is (in rad/s)
(a) 1.5p (b) 0.33p (c) 3p (d) 0.5p (e) 2p
Q.3 If you are given the force constant of a spring as 100 N/m, the mass on the spring as 1 kg, and the
amplitude as 0.04 m, then the period in seconds of the SHM is
(a) 20p (b) p/5 (c) 5p (d) 5/p (e) p/20
Q.4 If two vibrations of the same frequency are superimposed on a system with equal amplitudes, the
maximum resultant amplitude shall be
(a) zero (b) ÷(2A) (c) A/÷2 (d) 2A (e) A
Q.5 Speed of sound waves in a fluid
(a) depends directly on the density of the medium
(b) varies as square of bulk modulus of the medium
(c) is inversely proportional to the square root of density
(d) depends directly on the square root of the bulk modulus of the medium
Q.6 A transverse harmonic wave on a string is described by y(x, t) = 3.0 sin (36t + 0.018x + p/4) where x
and y are in cm and t is in s. The positive direction of x is from left to right. Then
(a) the wave is travelling from right to left (b) the speed of the wave is 20 m/s
(c) the frequency of the wave is 5.7 Hz (d) the least distance between two successive crests in
the wave is 2.5 cm
Q.7 An object is vibrating at its natural frequency. Repeated and periodic vibrations of the same natural
frequency impinge upon the vibrating object and the amplitude of its vibrations are observed to
increase. This phenomenon is known as
(a) beats (b) resonance (c) interference (d) overtone
Q.8 Standing waves are produced in a wire by vibrating one end at a frequency of 100 Hz. The distance
between the 2nd and the 5th nodes is 60.0 cm. The wavelength of the original traveling wave in cm is
(a) 50.0 (b) 40.0 (c) 30.0 (d) 20.0
Q.9 Which phenomena can be applied to estimate the velocity of star with respect to earth
(a) Doppler effect (b) interference of waves
(c) beats phenomena (d) all of these
Q.10 Doppler Effect applies to
(a) sound waves only (b) light waves only
(c) both sound and light waves (d) neither sound wave nor light waves
Q.11 The Lissajous patterns help in the measurement of
(a) Phase difference between two sine waves
(b) Frequency of one waveform if the frequency of other waveform is known
(c) both (a) and (b)
(d) none of these
Q.12 If the two input waveforms of equal amplitude and 90 degree phase difference is applied to the CRO,
then the Lissajous patterns obtained will be
(a) straight line tilted at 45 degree with respect to x-axis
(b) circle
(c) ellipse
(d) vertical straight line
Q.13 A body executing SHM has a velocity of 2.0 cm/s when its displacement is 7.0 cm and a velocity of
7.0 cm/s. What is the square of the amplitude of oscillation when its displacement is 2.0 cm?
(a) 26.0 cm2 (b) 53.0 cm2 (c) 79.0 cm2 (d) 106.0 cm2
Q.14 A body executing linear SHM has a velocity of 3.0 cm/sec when its displacement is 8.0 cm, and
a velocity of 8.0 cm/sec when its displacement is 3.0 cm. If the oscillator mass is 5.0 kg, find the
approximate total energy of the oscillator?
(a) 4.5 mJ (b) 9.1 mJ (c) 13.7 mJ (d) 18.2 mJ
Q.15 Two block, each of mass m = 2.0 kg, are connected by a spring of force constant R = 3.0 N/m and
placed on a horizontal frictionless surface, as shown in the following diagram. If an equal force of
F = 2.0 N is applied to each block in the direction of arrow, what is the approximate time-period of the
system when the force is removed.
(a) 1.2 sec (b) 2.4 sec (c) 3.6 sec (d) 4.8 sec
Q.16 In SHM of a simple pendulum, component of weight directed towards mean position is
(a) mg cos q (b) mg sin q (c) 0 (d) mg tan q
Q.17 Which of the following quantities are always positive in a SHM?
       
(a) F , a (b) V , r (c) a, r (d) F , r
Q.18 A small block oscillates back and forth on a smooth concave surface of radius R. The time period of
small oscillation is
2R
(a) T = 2p R /g (b) T = 2p (c) T = 2p R /2g (d) None of these
g
Q.19 When two mutually perpendicular simple Harmonic motions of same frequency, amplitude and phase
are superimposed,
(a) Resulting motion is Uniform circular motion.
(b) Resulting motion is a linear SHM along a straight line inclined equally to the straight lines of
motion of component ones.
(c) Resulting motion is an elliptical motion, symmetrical about the lines of motion of the component.
(d) The two SHM’s will cancel each other.
Q.20 A simple pendulum has some time period T. What will be the percentage change in its time period if
its amplitude is decreased by 5%?
(a) 6% (b) 3% (c) 1.5% (d) 0%
P ractice P roblems
General Questions
Q.1 On what factors the velocity of the sound depends?
Q.2 Do you agree that sound waves are mechanical waves?
Q.3 How are stationary waves formed?
Q.4 Define simple harmonic motion (SHM). Give two examples of SHM. Why SHM is important in the
study of waves and oscillations?
Q.5 Name three parameters of any SHM. Explain the meaning of each of them.
Q.6 Derive an expression for the energy of a harmonic oscillator of mass m, amplitude A, and frequency v.
Find out the displacement at which energy is half kinetic and half potential.
Q.7 All simple harmonic motions are periodic but all the periodic motions are not simple harmonic.
Explain this observation.
Q.8 Briefly explain how a simple harmonic motion can be represented by a rotating vector.
What is a shock waves?
Q.9 How are sound waves different from shock waves?
Q.10 Define the intensity of the sound?
Q.11 What do you understand by beats? Discuss their theory.
Q.1 The force constant of a spring is 10 N/m. Find the period of a 100-g mass on the end of this spring.
Ans. 0.63 s
Q.2 Find the maximum velocity of the mass in problem 1 if the amplitude of oscillation is 2.0 cm.
Ans. 20 cm/s
Q.3 Find the velocity of the mass in Problem 2 when it is 1 cm from its equilibrium position.
Ans. 17 cm/s
Q.4 A mass on the end of a spring is released from a point 2 cm from its equilibrium position. The frequency
of oscillation is 4 Hz. Write the equation for the position of the mass as a function of time.
Ans. 0.02 cos (8pt) m
Q.5 The period of a simple pendulum is 2.00 sec, find the length of the pendulum. Ans. 0.993 m
Q.6 Find the maximum energy stored in the spring of Problem 1 when it is compressed 2 cm from its
equilibrium position. Ans. 0.02 J
Q.7 The speed of sound waves in air is found to be 340 m/s. Determine the fundamental frequency (1st
harmonic) of an open-end air column that has a length of 67.5 cm. Ans. 252 Hz
Sound Waves and

Acoustics of Buildings
9
Learning Objectives
L01 Know about audible, ultrasonic, and L04 Analyse the applications of ultrasonic
infrasonic waves waves
L02 Learn about Production of ultrasonic L05 Explain the types of acoustics
waves, magnetostriction method, L06 Discuss on acoustics of buildings
Piezoelectric method, ultrasonic
L07 Evaluate the factors affecting the
transducer, Galton Whistle method
architectural acoustics
L03 Understand the absorption, dispersion,
and detection of ultrasonic waves
Introduction
A vibration refers to the oscillating motion of any medium and sound is a vibration in an elastic medium.
These vibrations transmitting through a solid, liquid, or gas, are composed of frequencies within the range
of hearing and are of a level sufficiently strong to be heard. In the case of human hearing it is the vibrations
in air that simulate our hearing organs and give a sensation of sound. When sound enters a new medium,
it is reflected, transmitted, or absorbed. This scientific study of the propagation, absorption, and reflection
of sound waves is called acoustics.
Acoustics is the interdisciplinary science that deals with the study of sound, ultrasound and infrasound (all
mechanical waves in gases, liquids, and solids). In a broad sense, acoustics may be defined as generation,
transmission and reception of energy in the form of vibration waves in matter.
9.1 AudibLe, uLtrAsonic And infrAsonic WAves LO1
The simplest form of sound waves is sinusoidal waves of definite frequency, wavelength and amplitude. The
frequency range of waves from 20 Hz to 20,000 Hz are said to be audible waves for which range human ears
are sensitive but the waves of frequency above the audible range are called ultrasonic waves and below the
audible range are known as infrasonic waves.
Sound Waves and Acoustics of Buildings 285
Ultrasonics is the study and application of the energy of sound waves vibrating at frequencies greater than
20,000 Hz, i.e., beyond the range of human hearing. The application of sound energy in the audible range
is limited almost entirely to communications, since increasing the pressure, or intensity, of sound waves
increases loudness and therefore causes discomfort to human beings. Ultrasonic waves, however, being
inaudible, have little or no effect on the ear even at high intensities. They are produced, commonly, by a
transducer containing a piezoelectric substance, e.g. a quartz-crystal oscillator that converts high-frequency
electric current into vibrating ultrasonic waves.
Sound waves, particularly in the atmosphere, whose frequencies are below the audible range, i.e., lower than
about 20 Hz are called infrasonic waves. Earthquake and seismic waves are elastic waves which occur at
infrasonic frequencies in the Earth’s crust and in the oceans and seas. The physical laws of propagation in
the atmosphere are essentially the same as for audible sound. The local speed of infrasound in air at ambient
temperatures near 20°C is about 340 m/s, the same as for audible sound.
9.2 Production of uLtrAsonic WAves LO2

In most applications, ultrasonic waves are generated by applying an electric current to a special kind of
crystal known as a piezoelectric crystal. The crystal converts electrical energy into mechanical energy, which,
in turn, causes the crystal to vibrate at a high frequency. In another technique, a magnetic field is applied to
a crystal, causing it to emit ultrasonic waves. Although the bulk attention is given to more popular types of
transducers which are based on magnetostriction and piezoelectric effect, there are other means of generating
ultrasonic waves some of which exhibit great promise.
9.2.1 Magnetostriction Method
Before discussing this method for the generation of ultrasonic waves, we shall talk about the magnetostriction
effect.
9.2.1.1 Magnetostriction Effect
When a rod of ferromagnetic material such as iron, nickel or cobalt is placed in a magnetic field keeping
its length parallel to the direction of magnetic field, the rod experiences a small change in its length. This
effect is termed as magnetostriction effect. The change in length of the rod depends on the intensity of
the applied magnetic field and nature of the ferromagnetic material. However, the change in the length is
independent of the direction of the field. Since the change is not so great in the other dimensions of the rod,
the rod is generally put with its length parallel to the direction of the magnetic field. The cause of change in
material’s dimensions can be understood as follows. Actually ferromagnetic materials have a structure that is
divided into domains, each of which is a region of uniform magnetic polarisation. Under the application of an
external magnetic field, the boundaries between the domains shift and the domains rotate. These two effects
lead to a change in the dimensions of the materials.
9.2.1.2 Principle Involved
The general principle involved in producing ultrasonic waves is to cause ferromagnetic materials to vibrate
very rapidly. These vibrations cause surrounding air to vibrate with the same frequency, which spreads out in
the form of ultrasonic waves.
When the rod is placed inside a magnetic coil carrying alternating current, it suffers a change in length for
each half of the alternating current. It means the rod vibrates at a frequency twice that of the frequency of the
alternating current. Usually the amplitudes of vibrations are small, but these can be enhanced by achieving
the resonance condition, i.e., by matching the

frequency of the alternating current with the
natural frequency of the material of the rod.
1 2
9.2.1.3 Construction and Working
In Fig. 9.1, a rod (ferromagnetic material) with
its ends A and B is winded by the coils L1 and
L2. The coil L2 is connected to the collector of
the transistor whereas the coil L1 is connected
to its base. In view of an LC circuit, we can

adjust the frequency of the oscillatory circuit
Ê 1 ˆ
ÁË 2p L C ˜¯
by adjusting the value of the Figure 9.1
2
capacitor C. The current flowing in the circuit
can be determined by the milliammeter connected across the coil L2. A necessary biasing, i.e., the emitter as
forward biased and the collector as reverse biased for the NPN transistor, is achieved by the battery (current)
connected between the emitter and the collector of the transistor. The alternating current passing through the
coil L2 causes a corresponding change in the magnetization of the rod and hence the rod starts vibrating due
to the magnetostriction effect.
In the above situation, an emf is also induced in the coil which is called as converse magnetostriction effect.
Due to this effect an emf is induced in the coil L1. This induced emf is fed to the base of the transistor, which
acts as a feed back continuously. This way the current is built up in the transistor and the vibrations of the rod
are maintained for the generation of ultrasonic waves. When the frequency of the oscillatory circuit matches
with the natural frequency of the vibrating rod, the resonance occurs. At the resonance, the rod vibrates
longitudinally with larger amplitude and produces ultrasonic waves of high frequency along both the ends of
the rod.
If the Young’s modulus of the material of the rod is Y, its density is r and the length of the rod is l, then the
1 Y
frequency of vibrations of the rod is given by . When this frequency matches with the frequency of the
2l r
1
oscillatory circuit , the resonance occurs. Based on this we get
2p L2C
1 1 Y
= .
2p L2C 2l r
9.2.1.4 Advantages and Limitations

There are several advantages of magnetostriction method of generating ultrasonic waves. For example,
magnetostrictive or ferromagnetic materials are easily available and inexpensive. The design of the oscillatory
circuit is simple and it involves low cost of the materials. Moreover, at low ultrasonic frequencies, the large
power output can be produced without any risk of damage of the oscillatory circuit.
On the other hand, this method has some limitations also. For example, through this method we cannot
achieve ultrasonic waves of frequencies larger than about 3 MHz. The frequency of oscillations also depends
on the temperature and degree of magnetization. So hearing effect may change the frequency of ultrasonic
waves and there will be losses of energy due to hystersis and eddy current. Finally the condition of resonance
shows that we need to reduce the length of the rod in order to produce higher frequency ultrasonic waves,
which is not practically feasible.
9.2.2 Piezoelectric Method LO2

Before discussing this method for the generation of ultrasonic waves, we shall talk about the piezoelectric
effect.
9.2.2.1 Piezoelectric Effect
In 1880, the brothers Pierre and Jacques Curie discovered in an experiment the fact that certain crystals
can develop an electric charge when a mechanical pressure or tension is applied. This phenomenon was
later named as Piezoelectric effect. They showed that there was a direct proportion between the mechanical
pressure and the resultant charge, and sign of the charge changed when pressure changed to tension or vice
versa. This may take the form of a separation of electric charge across the crystal lattice. If the material is not
short-circuited, the applied charge induces a voltage across the material. The crystals which acquire a charge
when compressed, twisted or distorted are said to be piezoelectric. The effect is present in many crystals but
it is useful in Quartz and 6.
In addition, certain ceramics and biological matter such as DNA, bone and various proteins show the
piezoelectric effect or the piezoelectricity in response to applied mechanical stress. The word piezoelectricity,
which means electricity resulting from pressure, was derived from the Greek words piezō or piezein that
means to squeeze or press and ēlektron that means amber (an ancient source of electric charge). Piezoelectric
effect finds useful applications such as the production and detection of sound, generation of high voltages,
electronic frequency generation, microbalances, and ultrasonic focusing of optical assemblies.
Quartz
1 2
Milliammeter
Key
NPN
Battery
Figure 9.2
When a slab of a piezoelectric crystal such as quartz is placed between two metal plates and resonant
mechanical vibrations are produced in the crystal due to the linear expansion and contraction, elastic waves
are propagated in the metallic plates which generate ultrasonic waves. An efficient generation of ultrasonic
waves takes place when the crystal oscillates at the maximum amplitude. This happens when the frequency
of the oscillatory circuit matches with the natural frequency of one of the modes of vibrations of the crystal.
The frequency of the generated ultrasonic waves depends on the Young’s modulus and the density of the
piezoelectric material.
Piezoelectric generator that works on the piezoelectric effect is used for generating ultrasonic waves of high
frequency of about 50 MHz. For this a slice of quartz crystal is placed between two metal plates A and B in
order to form a parallel plate capacitor having the quartz crystal as a dielectric medium. Quartz is preferred
because it possesses rare physical and chemical properties. The metal plates are connected to the terminals
of a coil which is inductively coupled to the oscillating circuit, as shown in Figure 9.2. Due to this electrical
circuit, an alternating potential difference is developed across the plates of the capacitor because of which a
tensile pressure appears on the crystal. This produces alternate contraction and expansion of the crystal and
the opposite charges are generated on the faces of the crystal lying towards A and B. Through piezoelectric
effect the crystal produces sound waves and when the frequency of electrical oscillations is in the ultrasonic
range then ultrasonic waves are generated.
As shown in Fig. 9.2, the variable capacitor C is adjusted in order to match the frequency of the oscillatory
circuit with the natural frequency of one of the modes of vibrations of the crystal. This way we are able
to produce resonant mechanical vibrations in the crystal due to the linear expansion and contraction. If
one or both the faces of the crystal are placed in contact with some medium in which elastic waves can be
propagated, ultrasonic waves are generated. The LC circuit having a variable capacitor C and an inductor L2
decides the frequency of the electrical oscillations. When the circuit is closed, the current flows through the
LC circuit and the capacitor is charged. The current stops flowing when the capacitor is fully charged. After
that the capacitor is made to discharge through the inductor so that the electric energy is stored in the form
of electric and magnetic fields associated with the capacitor and the inductor, respectively. This way we get
electrical oscillations in the circuit and with the help of the other electronic components including a transistor,
electrical oscillations are produced continuously. This is fed to the secondary circuit and the crystal vibrates,
as it is continuously subjected to alternating electric field.
9.2.3 ultrasonic transducer LO2

A transducer is a device which is used to convert one form of energy to another. Ultrasonic transducers convert
electrical energy to mechanical energy and vice versa. Ultrasonic sound can be produced by transducers which
operate either by the piezoelectric effect or the magnetostrictive effect. The magnetostrictive transducers can
be used to produce high intensity ultrasonic sound in the 20–40 kHz range for ultrasonic cleaning and other
mechanical applications. Ultrasonic transducers are constructed by incorporating one or more piezoelectric
vibrators which are electrically connected to pulsing-receiving system. An ultrasonic transducer includes
an ultrasonic transmitting/receiving element typically consisting of piezoelectric element connected to
electrodes. The piezoelectric elements typically are made of material such as lead zirconate titanate (PZT),
with a plurality of elements being arranged to form a transducer assembly. The transducer assembly is then
further assembled into a housing possibly including control electronics, in the form of electronic circuit
boards, the combination of which forms an ultrasonic probe.
The active element is the heart of the transducer as it converts the electrical energy to acoustic energy, and
vice versa. The active element is basically a piece of polarised material (i.e., some parts of the molecule are
positively charged, while other parts of the molecule are negatively charged) with electrodes attached to two
of its opposite faces. When an electric field is applied across the material, the polarised molecules will align
themselves with the electric field, resulting in induced dipoles within the molecular or crystal structure of
the material. This alignment of molecules will cause the material to change dimensions. This phenomenon
is known as electrostriction. In addition, a permanently-polarised material such as quartz (SiO2) or barium
titanate (BaTiO3) will produce an electric field when the material changes dimensions as a result of an
imposed mechanical force.
The thickness of the active element is determined by the desired frequency of the transducer. A thin wafer
element vibrates with a wavelength that is twice its thickness. Therefore, piezoelectric crystals are cut to a
thickness that is half the desired radiated wavelength. The higher the frequency of the transducer, the thinner
is the active element. The primary reason that high frequency contact transducers are not produced is because
the element is very thin and too fragile.
9.2.3.1 Uses of Ultrasonic Transducers
Ultrasonic transducers are useful for various applications. Ultrasonic testing equipment is used in a variety
of applications such as for measuring flow, determining flaws, measuring thickness, and gauging corrosion.
Ultrasonic diagnostic imaging systems are in widespread use for performing ultrasonic imaging and
measurements of the human body through the use of probes which are used to view the internal structure of
a body by creating a scan plane.
9.2.4 Galton Whistle Method LO2

Galton’s whistle was invented in 1876 by Francis Galton and is mentioned in his book named Inquiries into
Human Faculty and its Development. In this book, he described experiments to test the range of frequencies
that could be heard by various animals. This whistle most commonly known as dog whistle or silent whistle
is a type of whistle that emits sound in the ultrasonic range, which people cannot hear but some other animals
can, including dogs and domestic cats.
Galton and subsequent researchers used these whistles to create increasingly higher frequency tones to test
research subjects as well as animal abilities to hear different tones, Galton was able to determine that the
normal upper limit of human hearing was about 18 kHz. He also noted that the ability to hear higher
frequencies declined with age.
Below we discuss the principle involved and the construction and working of the
Galton’s whistle with regard to the production of ultrasonic waves. N
Galton whistle works on the principle of organ pipe, where the distance of annular
nozzle from the edge of a pipe and the pressure of air blast are suitably adjusted
in order to set the pipe into resonant vibrations at the ultrasonic frequency with
the help of the length and the diameter of the pipe.
As shown in Fig. 9.3, Galton whistle consists of a closed end air column whose
length can be adjusted with the help of a movable piston P. A screw S is connected
to this piston which can move the piston to the desired position. The open end Figure 9.3
of the pipe O is fitted with a lip L, and the gap between the ends O and A can be adjusted with the help of
another screw SN which can move the pipe A up or down. A nozzle N is fitted on the top through which an
air blast is blow towards lip L. When the blast of air strikes against the lip L, the column of air in the pipe is
set into vibration. The resonant position is achieved in order to produce the ultrasonic waves by adjusting the
length of the air column in O. Clearly the resonance frequency depends on the size of the pipe, i.e., its length
and diameter.
The wavelength l of the sound wave depends on the length l of the air column in O and the end correction x.
This is given by
l = 4(L + x)
From this we can calculate the frequency of the sound or ultrasonic wave as
V V
f = =
l 4(l + x)
Here V is the velocity of the waves produced by Galton’s whistle. This whistle can produce ultrasonic waves
of low frequencies up to 100 kHz and interestingly the micrometer screw S can be calibrated to give directly
this frequency.
9.3 AbsorPtion And disPersion of uLtrAsonic WAves LO3

When an ultrasonic wave passes through a medium, a part of its energy is converted into heat due to
the alternative compression and rarefaction taken place in the wave phenomenon and hence its intensity
goes down. The compressions produce the heat that increases the temperature of the medium whereas the
rarefactions reduce the temperature, leading to the absorption of these waves in the medium and the wave is
said to be attenuated. Two main mechanisms namely absorption and scattering (dispersion) are responsible
for the ultrasound attenuation. Different mechanisms such as thermal conductance effects, chemical effects,
viscous effects and nonlinearity are responsible for the absorption phenomenon. The phenomena responsible
for the ultrasound absorption in biological tissues have not been so far completely understood. In liquids
the viscous forces between neighbouring particles moving with different velocities are the major sources
of the wave absorption, whereas viscoelastic forces are the main contributors to the wave absorption in
homogeneous solids. For example, viscous losses may explain well the sound wave absorption in water
where attenuation varies with the square of the frequency. However, this model of viscosity does not explain
the experimental measurements of the absorption of the ultrasonic waves in soft biological tissues and bone in
the diagnostic frequency range. We can say that the absorption and dispersion of ultrasonic waves generally
focus on the influence of ultrasonics on molecular processes in liquids and gases, including hydrodynamics,
energy exchange, and chemical reactions.
The diminution in intensity of the amplitude of a planar longitudinal wave passing through a liquid is caused
by the conversion of the organised collective motion of the sonic pulse into random thermal motion, if we
neglect the radiation losses. The total attenuation comprises contributions arising from viscous loss and
thermal conductivity. In general, overall attenuation of the ultrasonic waves is characterised by the following
exponential decrease of the pressure amplitude p and of the intensity amplitude I with the propagating
distance z.
p = p0e–az and I = I0e2–a z
Here p0 and I0 are the pressure and intensity at z = 0, i.e., when the wave starts penetrating the medium. The
quantity a is called the pressure frequency-dependent attenuation coefficient, which is expressed in cm–1. The
factor of 2 in the exponential term of the intensity equation results from the transformation of the pressure
into intensity, as the intensity is proportional to the square of the pressure. The commonly used units for a in
biomedical ultrasonics are dB (decibel).
The dispersion of the ultrasonic wave is referred to the change in its velocity with frequency. In viscous liquids
such as glycerine and castor oil the change in velocity with frequency or dispersion cannot be observed in the
frequency regime or ultrasonic waves. However, the dispersion of these waves has been observed indirectly
by determining the change in wavelength of the waves.
9.4 detection of uLtrAsonic WAves LO3

There are several methods of detecting the ultrasonic waves, which include Kundt’s tube method, sensitive
method, and piezoelectric detection method.
Heaps of powder
Rod clamped at the centre
l
l Wave in air
2
Wave in rod
Figure 9.4
9.4.1 Kundt’s tube Method

Kundt’s tube has been very efficiently used for the detection of ordinary sound waves and the similar method
can be employed for detecting the ultrasonic waves. This tube is a long glass tube supported horizontally
with an air column in it. A horizontal rod is clamped at the center of the tube, as shown in Fig. 9.4. This
tube contains lycopodium powder scattered in it. When the ultrasonic waves are passed through the tube, the
lycopodium powder collects in the form of heaps which are found to be situated at the nodal points whereas
the powder is found to be blown off at the antinodal points. The average distance between two adjacent heaps
gives rise to the value of half wavelength from which the wavelength of the waves can be calculated. In
view of the use of the powder, the method is suitable for the detection of the ultrasonic waves of appreciable
wavelengths and it cannot be employed if the wavelength of the waves is very small, i.e., less than few
millimeters. However, in the case of liquid medium, powered coke is used in place lycopodium powder to
detect the position of nodes and the wavelength of the waves.
If the average distance between the adjacent heaps is d and the frequency of the ultrasonic waves is f, then the
velocity of the wave is given by V = 2fd.
9.4.2 sensitive flame Method

This method works on the basis of interaction of wave and a sensitive flame, where the change in pressure
is noticed. In this method of detection of ultrasonic waves, a narrow sensitive flame is moved along the
medium and the change in its intensity is noticed. At the positions of antinodes, the flame is found to be
steady (stationary), while the flame is found to flicker at the positions of nodes due to a change taken place in
the pressure. The positions of the nodes and antinodes are found out in the medium, and the average distance
between the two adjacent nodes gives rise to the value half wavelength. If the value of the frequency of
ultrasonic wave is known, the velocity of the wave passing through the medium can be calculated using the
same formula as used in the method of Kundt’s tube.
9.4.3 Piezoelectric detector
Piezoelectric effect, which is being used in the production of ultrasonic waves based on quartz crystal, can also
be used to detect the ultrasonic waves. The underlying principle is as follows. If ultrasonic waves comprising
of compressions and rarefactions are allowed to fall upon a quartz crystal, a certain potential difference is
developed across the faces of the crystal and varying electric charges are produced. These small charges after
amplification by an electronic circuit are used to detect the ultrasonic waves.
9.5 APPLicAtions of uLtrAsonic WAves LO4

Ultrasonics have found diverse applications in various fields of medical science or medicine, industry and
communication.
9.5.1 Medical Applications
After the discovery of X-ray imaging in the late 19th century, great advances have been made to diagnosis and
treatment equipment based on ultrasonics.
9.5.1.1 Diagnosis
Scanning of internal organs, vessels and tissues of patient’s body based on ultrasonic waves is called
ultrasonography. This makes use of high frequency sound waves to produce the images of internal organs and
structures for the medical examination and it is possibly the best of all ultrasonic medical applications. The
ultrasonic scans are less costly, quicker and easier to use than MRI (magnetic resonance imaging) and CT
(computerized tomography) scans. Hence, these are frequently used to monitor and diagnose the condition of
organs such as kidneys, liver or gallbladder. In order to diagnose and follow up heart conditions, doctors make
efficient use of EVG (echocardiograms) or ultrasonic scans of the heart of the patient.
9.5.1.2 Surgery
The technology based on ultrasound is increasingly being used in surgery. Here ultrasonic surgical instruments
convert an ultrasonic signal into a mechanical vibration by using a transducer. A waveguide is then used to
amplify and propagate the vibration to a desired position. The ultrasonic surgical instruments are highly
useful in diverse medical procedures, as these can cut bone and other tissue. At the same time reduce bleeding
by coagulating tissue. Finally this reduces the average length of surgery and damage to tissue, resulting in
fewer complications only.
9.5.1.3 Non-invasive Therapeutic Applications
Ultrasound energy can be used as non- or minimally invasive high intensity focused ultrasound (HIFU) or
high intensity therapeutic ultrasound (HITU). By applying ultrasound energy to heat and destroy diseased
tissues, these methods can be used to remove body tissue while treating the cancers and other conditions.
Ultrasound imaging systems locate and target liver, kidney or gallbladder stones. These are smashed into
pieces by ultrasound pulses and are finally evacuated naturally through urination. Other treatments using
ultrasound technology include bone healing and physiotherapy for inflammation caused by joint injuries.
Drug delivery is also done based on HIFU/HITU to treat tumours, especially in the brain where it may be
difficult to achieve. Cosmetic applications, such as non-invasive liposuction and for a number of therapies to
improve skin tone, scars and sun based damage also make use of ultrasound technology.
9.5.1.4 Dental Care
Another application of ultrasonics is in dental care as descalers to remove plaque. Ultrasonic descalers have
a tip that vibrates at high frequency to break down the bacterial matter to which plaque and calculus stick.
The ultrasonic waves have been found quite useful for painless dental cutting. This technology enables a
smoother and less painful experience.
9.5.1.5 Hygiene Safely
All medical and dental equipment must be absolutely clean before use, otherwise the introduction of
pathogenic microbes can lead to infection. It is very important to clean, disinfect and sterlize all multiple
use instruments and devices after their use on a patient or surgery. In this direction, ultrasonic cleaning uses
a special wash solution to reach and effectively remove organic waste from difficult-to-clean areas, such as
equipment or devices with joints and crevices.
9.5.2 Industrial Applications
Industrial Applications of ultrasonics include ultrasonic machining, welding, cleaning, etc.
9.5.2.1 Machining
Ultrasonic machining is a vibratory process which is now in common use for the mechanical treatment of
hard and brittle solids such as glasses, ceramics, precious stones, semiconductors and hard alloys. A glass rod
oscillating with ultrasonic frequency can be used to bore holes in steel and other hard metals.
9.5.2.2 Welding
With regard to the application of ultrasonics for welding it is believed that practically all metals and plastics
can be welded ultrasonic waves of suitable energy. Here the ultrasonic energy converts into heat at the
contact area as a result of friction arising between the surfaces. As the temperature of surfaces’ layers exceeds
the crystallization point, both the layers melt and make a bond together to form a strong joint. Since this
process induces negligible stress at the spot of welding, this is quite attractive that the structure of materials
remains unchanged.
9.5.2.3 Cleaning
Towards the cleaning applications of ultrasound waves, it is worth mentioning that these waves with
frequencies 20 kHz to 40 kHz are used for cleaning of jewellery, optical parts, surgical instruments, industrial
parts etc. They are used for cleaning clothes and parts of watches. Printing industry used ultrasonic as a
method of cleaning complicated and problematic parts has been available for many years with in a wide
range of industries. The main advantages are that components of the most complicated shapes can be cleaned
efficiently, speedily and comprehensively. Here ultrasonic millions of tiny bubbles within the fluid which act
on the surface of the component behave as a brush in many ways. The scrubbing action of this brush can be
made as vigorous or gentle as per the requirement.
9.5.2.4 Structural Composition and Analysis
Ultrasonic waves are used for producing alloys of uniform composition. Further, these waves are employed
to detect cracks or flaws in metal structure.
9.5.3 Applications in food technology

By tuning frequency, ultrasound can be utilized in food technology. Since ultrasound techniques are relatively
cheap, simple and energy saving, these have became an emerging technology for probing and modifying food
products. Low power (high frequency) ultrasound is used for monitoring the composition and physiochemical
properties of food components and products during processing and storage. However, high power (low
frequency) ultrasound induces mechanical, physical and chemical (biochemical) changes through cavitation,
which supports many food processing operations such as extraction, freezing, drying, emulsification and
inactivation of pathogenic bacteria on food contact surfaces.
Using ultrasound, full reproducible food processes can now be completed in seconds or minutes. This can
be done with high reproducibility, reducing the processing cost, simplifying manipulation and giving higher
purity of the final product. This also eliminates post-treatment of waste water and consumes only a fraction
of the time and energy normally needed for conventional processes.
9.5.4 Applications in Communications
Ultrasonic waves can be produced in the form of beams in the desired direction. These can travel long
distances in water before being absorbed. This makes them suitable for the submarine applications. Submarine
ultrasonic transmitters have been developed for detecting the presence of iceberg or submarines. These are
used for signaling from ship to ship, especially in submerged submarines and also in determination of the
depth of sea, position of a ship and submarine. The ship is equipped with the source and receiver of a
particular frequency at its bottom. The source is used to transmit the short ultrasound pulses and the reflected
pulses are received by the receiver for the detection. Actually the time interval (t) between sending and
receiving the pulses is measured, which gives rise to the depth of the ocean as
Vt
d=
2
Here V is the velocity of the ultrasonic waves.
9.5.5 Detection of Velocity of Sound in Liquid

The stationary wave method is applied to find the velocity of ultrasonic waves in liquid and gases. The velocity
of ultrasonic waves in these medium can be calculated from the relation V = flult, where f is frequency and lult
is the wavelength of the ultrasonic waves in the medium. This method is more suitable for finding the velocity
of ultrasonics in liquid and gases that are available in small quantity.
In order to find the velocity of sound in a liquid, say kerosene oil, a quartz crystal is placed between two metal
plates and the plates are connected to an audio frequency oscillator. The assembly of the crystal and the plates
is kept inside the liquid cell. The crystal is made to vibrate in resonance with the oscillator by adjusting the
frequency to produce the ultrasonic waves. The reflections of the wave from the sides of the liquid cell form
a standing wave pattern with nodes and antinodes at regular intervals. This leads to a particular distribution
of the liquid density with maximum density at the nodes and minimum density at the antinodes. Accordingly
the refractive index of the liquid is varied and it works as the diffraction grating, known as acoustic grating.
The acoustic grating is mounted on a prism table of spectrometer and a parallel beam of sodium light from
collimator is allowed to fall normally on the grating. The diffracted light is found to form a diffraction pattern
which is viewed through a telescope. The diffraction pattern consists of central maximum and principal
maxima on either sides. The positions of the principal maxima satisfy the following relation
lult sin qn = nl
Here lult is the wavelength of sound in the liquid, l is the wavelength of incident sodium light (monochromatic)
and qn is the angle of nth other diffraction. We can find the wavelength lult of the wave in the liquid. If f be the
frequency of vibrations of the crysttal, then the velocity of the ultrasonic wave in the liquid can be obtained
using the relation V = flult.
9.6 tyPes of Acoustics LO5

There are various types of the acoustics. These are discussed below.
9.6.1 Physical Acoustics
Physical acoustics encompasses propagation and absorption of sound at all frequencies in air and other gases,
liquids, semi-solids and solids. It deals with airborne, audible sound, infrasound and ultrasound. Physical
acoustics includes both linear processes such as the propagation of sound from traffic, and nonlinear processes
such as the shock waves that are generated by planes flying faster than the speed of sound.
9.6.2 engineering Acoustics
Engineering acoustics deals with the development of devices to generate (e.g., loudspeakers), record (e.g.,
microphones) and analyse (e.g., frequency analysers) sound of all kinds. The field of sound production,
recording and reproduction, with all its attendant electronics and measuring instruments, is an important part
of engineering acoustics.
9.6.3 Architectural Acoustics
Architectural acoustics is concerned with sound in buildings. One aspect of this field is the control of sound
within rooms to maximise the acceptability of music or intelligibility of speech. This branch of architectural
acoustics deals with sound in lecture theatres, concert halls, meeting rooms and classrooms.
9.6.4 Musical Acoustics
Musical acoustics considers the workings of traditional, experimental and electronic musical instruments.
The interaction of musicians, instruments, listeners and performance spaces means that many branches of
acoustics influence work in this field.
9.6.5 Psychological Acoustics
Psychological acoustics studies the brain’s signal-processing function, which takes nerve impulses from the
ear and interprets them. Physiological acoustics deals with models and theories of the operation of the ear
and its anatomy. One practical application of this field is the study of the elements important to achieve a
stereophonic effect. Another is the determination of those factors that make one sound unpleasant or annoying
and another reverse. There is no direct correlation between loudness and annoyance.
9.6.6 bioacoustics
Bioacoustics studies all aspects of acoustic behaviour in animals and biological media in general. This field
includes topics such as sound production by animals, bio-sonar, sound reception by animals, effects of noise
on animals and medical diagnostics using acoustics, especially ultrasonics.
9.7 Acoustics of buiLdinGs LO6

The branch of the science which deals with the planning of a building or a hall with a view to provide best
audible sound to the audience is called acoustics of building or architectural acoustics. WC Sabine in 1911,
first of all scientifically tackled the problem of satisfactory speech and music in a hall.
9.7.1 reverberation
When a sound is produced in a building, it lasts too long after its production. It reaches to a listener a number
of times. Once it reaches directly from the source and subsequently after reflection from the walls, windows,
ceiling and floor of the hall. The listener, therefore, receives series of sounds of diminishing intensity (since
part of energy is lost at each reflection); the sound becomes muddy, garbled. The most important factor in the
design of an auditorium is reverberation. Reverberation is nothing but the prolonged reflection of sound from
the walls, floor and ceiling of a room. It is also defined as the persistence of audible sound after the source
has stopped to emit sound. The duration for which the sound persists is called reverberation time. The time of
reverberation is also defined as the time taken for the sound to fall below the minimum audibility measured
from the instant when the source stops sounding. Sabine, using an organ pipe of frequency 512 Hz found that
its sound becomes inaudible when its intensity fall to one millionth of its intensity just before stopping the
organ pipe. Hence, Sabine defines the standard revereberation time as the time taken by sound to fall to one
millionth of its intensity just before the source is cut off. Sabine found that the time of reverberation depends
upon the size of the hall, loudness of the sound and upon the kind of the music or sound for which hall is to
be used. For a sound of frequency 512 Hz, the best time of reverberation was found to be 1 to 1.5 sec and 1.5
to 2 sec for halls of 50,000 and 40,000 cubic feet, respectively.
Based on the range of values of revereberation time for specific purposes, we can determine a relationship
between room volume and internal surface area. This assumes the use of standard auditorium construction
materials.
9.7.2 Basic Requirement for Acoustically Good Halls

Before 1900, the architects and building engineers had no consideration about the acoustical properties of
rooms and halls etc. Sometimes, a building was found to be unsatisfactory for the purpose for which it was
built. According to Sabine the following essential features are required for the good acoustics.
(i) The sound heard must be sufficiently loud in every part of the hall and no echoes should be present.
(ii) The total quality of the speech and music must be unchanged, i.e., the relative intensities of the
several components of a complex sound must be maintained.
(iii) For the sake of clarity, the successive syllables spoken must be clear and distinct, i.e., there must be
no confusion due to overlapping of syllables.
(iv) The reverberation should be quite proper, i.e., neither too large nor too small. The reverberation time
should be 1 to 2 sec for music and 0.5 to 1 sec for speech.
(v) There should be no concentration of sound in any part of the hall.
(vi) The boundaries should be sufficiently sound proof to exclude extraneous noise.
(vii) There should be no Echelon effect.
(viii) There should be resonance within the building.
9.7.3 Transmission of Sound and transmission Loss

When sound is produced in a hall, it proceeds outward in spherical waves and strikes the boundaries of the
hall. The sound waves undergo reflection, absorption and transmission in varying amounts. The amounts of
these three processes depend upon the frequency of sound and the characteristics of the wall of the room (i.e.,
thickness of walls, weight, material, nature of surface). The transmission of sound means the sound energy
transmitted through the walls.
The loss of sound energy across a wall or a barrier is defined as transmission loss. So when the sound is
transmitted from the source to adjoining room or area, then there is a reduction in the sound intensity. This is
known as transmission loss. Thus, the transmission loss is numerically equivalent to loss in the intensity of
sound. This is expressed in decibels (dB).
(i) If 50 dB and 10 dB are the sound levels measured on either side of a wall, then transmission loss of
the wall would be, 50 dB – 10 dB = 40 dB.
(ii) Larger the transmission loss, greater is the sound illusion.
(iii) Transmission loss varies with material used for the construction.
(iv) The methods used for the construction also affect the value of transmission loss.
(v) Transmission loss varies with frequency of sound.
9.7.4 Acoustic environments and sound fields

There are two basic environments in which one makes measurements of sound and noise: outdoors and
indoors. An outdoor acoustic environment may be quite often referred to as a free field. A sound field is said
to be a free field if it is uniform, free from boundaries, and undisturbed by other sources of sound. Anechoic
chambers and well-above-the-ground outdoors are free fields. Sound radiated by a source in a free field
propagates away from the source and is never reflected back. Sound spreads three-dimensionally from the
source such that the intensity falls off as 1/r2 (6 dB decrease from source doubles).
The indoors acoustic environment introduces boundaries which reflect sound. If the boundaries completely
reflect all incident sound without any absorption then the resulting sound field is termed as diffuse
or reverberant. In a diffuse sound field the time average of the mean square sound pressure is the same
everywhere through the enclosure. The flow of energy is equally probable in all directions. If the boundaries
absorb some of the incident sound and reflect the rest, then the sound field is called semi-reverberant. Energy
flows in more than one direction. Semi-reverberant fields are the most widely encountered in the majority of
architectural acoustic environments.
9.7.5 indoor Acoustics
When a sound source is enclosed, the radiated sound energy is retained within the enclosure. If the boundaries
are perfectly reflective then the sound energy inside the enclosure could theoretically grow until a pressure is
reached that would be explosive. Fortunately, most realistic boundaries are at least partly absorbing (air also
absorbs sound) and the kinds of sound sources usually encountered in a room (for example, human speech)
are not extremely powerful. For example, the sound power produced by human speech is very small. Typical
male and female speakers generate 34 mW and 18 mW, respectively, at a distance of 3.28 ft. So, common
sound sources are not excessively powerful, the sound energy in the enclosure travels about the enclosure and
slowly decays as it is absorbed by the boundaries and the medium.
9.7.6 Sabine’s Formula for Reverberation Time

Consider a source of the sound in an enclosure. We shall assume that
(i) The rate of emission of energy from the source is constant.
(ii) The energy is transmitted equally in all directions and its distribution is uniform throughout the
enclosure.
(iii) The dissipation of energy in the air is negligible and that it is confined only to the boundary walls.
(iv) The power of absorption of the absorbent material is independent of the intensity of the sound.
(v) The effect of superposition is negligible.
The reverberation time can be measured by producing white noise in the enclosure at a level 60 dB above
background, switching it off, and measuring the time it takes to drop back to background. The drawbacks to
this method are that it takes a high power source and the noise produced is unpleasantly loud. Measuring the
time required for a smaller intensity drop and extrapolating to the 60 dB time is the usual procedure. This
requires some calculation and a mathematical model for the decay curve of the sound in the enclosure.
If a sound source of intensity I0 is switched off at time t = 0, then the sound starts to die away and will be
60 dB down (or 10–6 I0) after one reverberation time later, i.e., at t = T, where T is the revereberation time
(sometimes called RT60 as a mnemonic for its defintion). To evaluate the scaling constant a, note that
e–at/T = 10–6 at t = T (i)

where a is called scaling constant.
Taking the natural logarithm of both sides, we get
–a = –6 In (10) = –13.82
This evaluation will permit a calculation of the reverberation time T if any two intensities I0 and I are measured
along with the time interval t between the two measurements. From Eq. (i), the relationship for T is
13.28t
T= (ii)
ln ( I 0 /I )
It is found that rather than the intensity, the voltage is displayed from the microphone and the intensity is
proportional to the square of the voltage. In terms of the measured voltages, the reverberation time relationship
is obtained as
13.28t
T= (iii)
ln (V02 /V 2 )
This is the relationship which is to be used in the experiment to determine the reverberation time of the room.
It could then be compared with the approximate relationship
0.161 V
T= (iv)
aS
where V is room volume in cubic meters, S is surface area in square meters, and a is average absorption
coefficient. The rough absorption coefficient obtained in this manner could help to evaluate the change,
which might be made in the reverberation time by changing the surfaces (carpet, curtains, etc.).
Absorption coefficient for standard building materials is given by Sabine’s formula

0.161 V
T=
SaS
where SaS is the total acoustic absorption in the enclosure.
9.7.7 Absorption Coefficient and Its Measurement

What are the essential parameters of a typical room necessary to determine its acoustical behaviour? First,
an enclosed space that has an internal volume V. Second, it has a total boundary surface area S. Third,
each of the individual surface areas has an absorption coefficient. The average absorption coefficient for
all surfaces together is given by the individual surface areas and are the individual absorption coefficients
of the individual surface areas. Since the boundaries of the room reflect incident sound energy, the sound
signal received by a listener at some location in the room will consist of sound which arrives directly from
the source, sound which arrives after reflecting from one surface, and sound which has undergone several
reflections. The average distance between reflections in such a space is called the mean free path and is
related to the dimensions of the room. This is because each time the sound interacts with a surface in the
room it loses some of its energy due to absorption. The absorption coefficient of the material of a surface is
a, then a fraction (1 – a) of the sound energy incident upon the surface is returned to the side from which it
came. The absorption coefficient of sound is essentially the dissipation of energy into heat, and in so far as it
is affected by a bounding surface it is mainly due to one of two causes, porosity and flexural vibrations. The
absorption coefficient of any material, as originally defined by Sabine, is the ratio of the sound absorbed by
that material to that absorbed by an equivalent area of open window. Thus, a perfectly absorbent material
would have an absorption coefficient of 1.
Measurement Two audio frequency sources of powers P1 and P2 are employed in the test chamber. The
actual value of P1 and P2 need not be known, but their ratio must be known. Then
4V log( P1 / P2 )
A=
c T1 - T2
where T1 and T2 are the respective times of decay to the threshold of audibility. From this relation mean
coefficient of absorption can be obtained since A – aS, where S is the area of the absorbing surface is the
measure of mean coefficient of absorption. The test chamber for the experiment exclude all extraneous and
must have a long time of reverberation when empty.
Another reverberation method employed by Sabine was to measure the absorbing power of a room in terms
of that of an open window. Let aS and (aS + w) be the total absorption of the room with windows closed and
open, respectively and t0 and tc are the time taken for the sound to decay to the threshold of audibility in the
two cases, respectively. Then for same intensity of sound in both the cases,
w Ê t0 ˆ
= - 1˜
aS ÁË tc ¯
and for same energy P of source
Ê wˆ c
log Á1 + ˜ = (aS (tc - t0 ) - wt0 )
Ë aS ¯ 4V
From the above relation aS can be calculated in terms of w.
9.8 fActors AffectinG tHe ArcHitecturAL Acoustics LO7

By an acoustically good hall we mean that in which every syllable or musical note reaches an audible level of
loudness at every point of the hall and then quickly dies away to make room for the next syllable or group of
notes. The departure from this makes the hall defective acoustically. Following factors affect the architectural
acoustics:
(i) Reverberation: The reverberation can be controlled by the following factors:
(a) By providing windows and ventilators which can be opened and closed in order to optimise the
value of the time of reverberation optimum.
(b) By covering the floor with carpets.
(c) By heavy curtains with folds.
(d) By full capacity of audience.
(ii) Adequate Loudness: It can be reduced by using large sounding boards behind the speaker and
facing the audience. Large polished wooden reflecting surfaces immediately above the speakers are
also helpful. Low ceiling are also of great help in reflecting the sound energy towards the audience
by providing additional sound energy with the help of equipment like loudspeakers.
(iii) Focusing due to Walls and Ceiling: For uniform distribution of sound energy in the hall, there
should be curved surfaces. If such surfaces are present, they should be covered with absorbent
material. Ceiling should be low. A paraboloidal reflected surface arranged with the speaker at the
focus is also helpful in sending a uniform reflected beam of sound in the hall.
(iv) Echoes
(v) Resonance
(vi) Echelon Effect
(vii) Extraneous Noise and Sound Insulation: Generally, there are three types of noises, which are very
troublesome. These are Air Borne Noise, Structure Borne Noise and Inside Noise.
Air Borne Noise:
The noises are transmitted through the air. Sound insulation for the reduction of air borne noise can be
achieved by the following methods:
(a) By avoiding opening for pipes and ventilators.
(b) By allotting proper places for doors and windows.
(c) Using double doors and windows with separate frames and having insulating material between
them.
(d) By making arrangements for perfectly shutting doors and windows.
(e) Using heavy glass in doors, windows and ventilators.
(f) By providing double wall construction, floating floor construction, suspended ceiling construction,
box type construction, etc.
Structure Borne Noise:
The most common sources of this type of sound are footsteps, street traffic, hammering, drilling, operating
machinery, moving of furniture, etc. Sound insulation for the reducing of structure borne noise is done by
the following ways:
(a) By breaking the continuity by interposing layers of some acoustical insulator.

(b) By using double walls with air space between them.
(c) By using anti-vibrations mounts.
(d) By soft floor finishing (carpet, rubber, etc.)
(e) By insulating the machinery, as the mechanical equipments like refrigerators, lifts, fans etc., produce
vibrations in the structure.
Inside Noise:
The following methods are used for sound insulation of inside noise.
(a) The machinery (like typewriters) should be placed on absorbent pads.
(b) Any engine inside the hall should be fitted on the floor with a layer of wood.
(c) The floor should be covered with carpet.
(d) The wall, floors and ceiling should be provided with sound absorbing materials.
(e) The sound absorbing materials should be mounted on the surface near the source of noise.
9.8.1 sound Absorbing Materials

Special materials used to increase absorption of sound waves or to reduce the reflection of sound waves in a
room or hall are known as sound absorbing material. The material should have the following requirements:
(a) It should have good resistance to fire.
(b) It should be efficient over a wide range of frequencies.
(c) It should have high sound absorbing efficiency.
(d) It should be cheap, easily available, easy to fix, good looking, light in weight and waterproof
and should have economical maintenance and sufficient structural strength. The sound absorbing
materials are broadly classified into the following four categories, namely porous absorbents, cavity
resonators, composite types of absorbents, and resonant absorbents or panel absorbers.
s ummarY
✦ Scientific study of the propagation, absorption, and reflection of sound waves is called acoustics.
Acoustics is the interdisciplinary science that deals with the study of sound, ultrasound and infrasound
(all mechanical waves in gases, liquids, and solids). In a broad sense, acoustics may be defined as
generation, transmission and reception of energy in the form of vibration waves in matter.
✦ Various types of acoustics, namely physical acoustics, engineering acoustics, architectural acoustics,
musical acoustics, psychological acoustics, bioacoustics, were discussed.
✦ Description of audible waves, ultrasonic waves, and infrasonic waves were given.
✦ Certain crystals can develop an electric charge when a mechanical pressure or tension is applied. This
phenomenon is named as Piezoelectric effect.
✦ A transducer is a device which is used to convert one form of energy to another. Ultrasonic transducers
convert electrical energy to mechanical energy and vice versa. Ultrasonic sound can be produced by
transducers which operate either by the piezoelectric effect or by the magnetostrictive effect. The
magnetostrictive transducers can be used to produce high intensity ultrasonic sound in the 20–40 kHz
range for ultrasonic cleaning and other mechanical applications.
✦ Principle of ultrasonic transducer was discussed.
✦ It was discussed how ultrasonic waves are produced. Their applications were talked about.
✦ Acoustics of buildings was discussed in detail. Reverberation was introduced and it was said that the
reverberation is nothing but the prolonged reflection of sound from the walls, floor and ceiling of a
room. It is also defined as the persistence of audible sound after the source has stopped to emit sound.
The duration for which the sound persists is called reverberation time. The time of reverberation is also
defined as the time taken for the sound to fall below the minimum audibility measured from the instant
when the source stops sounding.
✦ Basic requirement for the acoustically good halls were discussed. These are the following.
(a) The sound heard must be sufficiently loud in every part of the hall and no echoes should be present.
(b) The total quality of the speech and music must be unchanged, i.e., the relative intensities of the
several components of a complex sound must be maintained.
(c) For the sake of clarity, the successive syllables spoken must be clear and distinct, i.e., there must
be no confusion due to overlapping of syllables.
(d) The reverberation should be quite proper, i.e., neither too large nor too small. The reverberation
time should be 1 to 2 seconds for music and 0.5 to 1 second for speech.
(e) There should be no concentration of sound in any part of the hall.
(f) The boundaries should be sufficiently sound proof to exclude extraneous noise.
(g) There should be no Echelon effect.
(h) There should be resonance within the building.
✦ Transmission of sound and transmission loss were discussed in detail.
✦ Sabine’s formula for reverberation time was derived and its theoretical as well as physical aspects were
talked about.
✦ Finally, the absorption coefficient was introduced and methods were talked about for its measurement.
✦ Factors affecting the architectural acoustics were discussed in detail and the methods of its removal
were talked about.
s olVeD e XamPles
E xamplE 1 The frequency limits of the range of human hearing ear is from about 20 Hz to 20 kHz. The
speed of sound is about 34,500 cm/sec. What is the wavelength of the wave in cm?
Solution The frequency range is given as 20 Hz to 20 kHz.
The speed of sound = 34500 cm/sec = 345 m/sec
We know that the frequency of sound wave is related to its wavelength by v = fl

345
\ l=
20
345
So for the frequency of 20 Hz, wavelength l = m = 17.25 m = 1725 cm
20
345
For the frequency of 20 kHz, wavelength l = m = 17.25 ¥ 10–3 m = 1.725 cm
20 ¥ 103
Therefore, wavelength range of the sound wave is 1.725 cm to 1725 cm.
E xamplE 2 Calculate the velocity of the sound in air in cm per sec at 100°C if the density of air at S.T.P. is
0.001293 g/cm3, the density of the mercury at 0°C is 13.60 g/cm3, the specific heat of air at constant pressure
is 0.2417 and the specific heat of air at constant volume is 0.1715.
Solution The velocity of sound in air is given by
gp
v= with usual notation.
r
The quantity r of the air is 0.001293 g/cm3

The pressure p is given by
p = hrg = 76 ¥ 13.6 ¥ 980 dynes/cm3
Cp 0.2417
Now g = =
Cv 0.1715
0.2417 ¥ 76 ¥ 13.6 ¥ 980

\ v=
0.1715 ¥ 0.001293
The velocity v is proportional to T , where T is the temperature of the air. Thus, if v¢ be the velocity at 100°C and v at
0°C, then
v¢ 273 + 100 373
= =
v 273 273
373
v¢ = v
273
373 ¥ 0.2417 ¥ 76 ¥ 13.6 ¥ 980
= = 38839.12 cm/sec
273 ¥ 0.1715 ¥ 0.001293
= 38839 cm/sec
E xamplE 3 The wavelength of the gas emitted by a tuning fork of frequency 512 vibration/sec in air at 17°C
is 66.5 cm. If the density of air at S.T.P. is 1.293 mg/cm3, calculate the ratio of two principal specific heats of
air. Assume that the density of mercury is 13.6 g/cm3.
Solution Since v = fl, the velocity of sound at 17°C is given by
v = 512 ¥ 66.5 cm per sec
gp
Now, v0 =
r
1.293
Here p = 76 cm of mercury = 76 ¥ 13.6¥ 980 dynes/cm3. The density of air r = g/cm3. If v0 be the velocity at 0°C
and since the velocity is proportional to T , 1000
273 273
v0 = v= ¥ 512 ¥ 66.5
290 290
gr
Now, v0 =
r
v02 r 273 ¥ (512 ¥ 66.5) 2 ¥ 1.293
\ g = = = 1.39
p 290 ¥ 1000 ¥ 76 ¥ 13.6 ¥ 980
E xamplE 4 A hall of floors is 15 ¥ 30 m2 along with height of 6 m, in which 500 people occupy upholstered
seat and the remainder sit on wooden chairs. Optimum reverberation time for orchestral music is 1.36 sec and
absorption coefficient per person is 0.44.
(a) Calculate the coefficient of absorption to be provided by the walls, floor and ceiling when the hall is
fully occupied.
(b) Calculate the reverberation time if only the half upholstered seats are occupied.
Solution
(a) The optimum reverberation time is T = 1.36 sec
Using Sabine’s formula equation of SI unit
V
T = 0.161
aS
0.161 ¥ (15 ¥ 30 ¥ 6)
1.36 =
aS
aS = 319 SI units
Absorption due to audience = 500 ¥ 0.44
= 220 SI units
Therefore, the absorption provided by the walls, floor and ceiling is
319 – 220 = 99 SI unit
(b) When the hall is only half filled the absorption will also be provided by vacant seats in addition to the absortion by
the audience.
250 ¥ 0.44 = 110 SI unit
The absorption by vacant wooden seats = 250 ¥ 0.02 = 5 SI unit
So the total absorption of the hall = 99 + 110 + 5 = 214 SI unit
Here the reverberation time, given by Sabine’s formula, is now
0.161 ¥ (15 ¥ 30 ¥ 6) 0.161 ¥ (15 ¥ 30 ¥ 6)
T= =
214 214
2.03 sec
E xamplE 5 Calculate the total absorption coefficient of cinema hall, whose volume is 8000 m3 and
reverberation time required is 1.8 sec.
Solution
The reverberation time is given by
0.161V 0.161V
T= =
aS Total absorption in hall
0.161V 0.161 ¥ 8000
\ Total absorption in hall = =
T 1.8
= 715.55 O.W.U.
E xamplE 6 Find out reverberation time of empty hall of volume 1700 m3 having a seating capacity for 150
persons with following data
Surface Area Coefficient of

absorption in O. W. U.
Plastered wall 98 m2 0.03
2
Plastered ceiling 144 m 0.04
2
Wooden door 15 m 0.06
2
Cushioned chairs 88 m 1.0
Solution Given V = 1700 m3

Based on the given data, the absorption by
Plastered wall = 98 ¥ 0.03 = 2.94
Plastered ceiling = 144 ¥ 0.04 = 5.76
Wooden door = 15 ¥ 0.06 = 0.90
Cushioned chairs = 88 ¥ 1.0 = 88.0
\ Total absorption = 97.6
0.161V 0.161 ¥ 1700
Reverberation time T = =
aS 97.6
or T = 2.80 sec
E xamplE 7 Calculate the reverberation time for a hall of volume 1400 m3, which has seating capacity of 110
persons with full capacity of audience and when audience are occupying only cushioned seats. Relevant data
may be taken from Ex. 6.
Solution We have total absorption in hall (from Ex. 6) = 97.6, V = 1400 m3.
When the hall is with full capacity of 110 person, the absorption due to them
= 110 ¥ 4.7 = 517
Now total absorption = 97.6 + 517 = 614.6
Reverberation time
0.161V 0.161 ¥ 1400
T= =
aS 614.6
T = 0.367 sec
E xamplE 8 The volume of a room is 980 m3. The wall area of the room is 150 m2, ceiling area is 95 m2 and
floor area is 90m2. The average sound absorption coefficient (i) for wall is 0.03, (ii) for ceiling is 0.80 and
(iii) for the floor is 0.06. Calculate the average sound absorption coefficient and the reverberation time.
Solution The average sound absorption coefficient

a1S1 + a2 S2 + a3S3
a=
S1 + S2 + S3
0.03 ¥ 150 + 0.80 ¥ 95 + 0.06 ¥ 90
=
150 + 95 + 90
= 0.256
and total area S = 150 + 95 + 90 = 335
Now total absorption of the room
= aS
= 0.256 ¥ 335 = 85.76 metric Sabines
Reverberation time
0.161V
T=
aS
0.161 ¥ 980
=
85.76
= 1.84 sec
E xamplE 9 How much acoustic power enters the window of area 1.58 m2, via the sound wave (standard
intensity level = 10–16 W/cm2). The window opens on a street where the street noise results in an intensity
level at the window of 60 dB.
Solution Given the intensity level at window = 60 dB
Area of the window = 1.58 m2
Standard intensity level I0 = 10–16 W/cm2 = 10–12 W/m2.
We know that intensity level = 10 log10(I/I0) dB
\ 60 = 10 log10 (I/10–12) dB
I = 9.98 ¥ 10–7 W/m2
Acoustic power = intensity ¥ area = 9.98 ¥ 10–7 ¥ 1.58 = 1.576 ¥ 10–6 W = 1.58 ¥ 10–6 W
E xamplE 10 Find the frequency to which a piezoelectric oscillator circuit should be turned so that a
piezoelectric crystal of 0.1 cm thickness vibrates in its fundamental mode to generate ultrasonic waves.
Young’s modulus and density of material of the crystal are 8 ¥ 1010 Nm–2 and 2.654 ¥ 103 kg m–3 respectively.
Solution Given, thickness of the crystal t = 1 ¥ 10–3 m, density (D) = 2.654 ¥ 103 and Y = 8 ¥ 1010 Nm–2
From the relation, the fundamental frequency of piezoelectric oscillator
p Y 1 8 ¥ 1010
f = =
2t D 2 ¥ 0.001 2.654 ¥ 103
2.75 106 Hz
E xamplE 11 Calculate the natural frequency of 30 mm of iron rod. The density of iron rod and Young’s
modulus are 7.25 ¥ 103 kg/m3 and 115 ¥ 109 N/m2 respectively. Can you use it in magnetostriction oscillator
to produce ultrasonic waves?
Solution Given, l = 3 ¥ 10–2 m, D = 7.25 ¥ 103 kg/m3 and Y = 115 ¥ 109 N/m2
From the relation, the frequency of ultrasonic waves by magnetostriction oscillator is
1 Y 1 115 ¥ 109
f = =
2l D 2 ¥ 3 ¥ 10-2 7.25 ¥ 103
= 66.38 ¥ 103 Hz
= 66.38 kHz
E xamplE 12 Calculate the fundamental frequency of a quartz crystal of 3 ¥ 10–3 m thickness. The density of
the crystal is 2650 kg m–3 and Young’s modulus is 7.9 ¥ 1010 N/m2.
Solution Given, l = 3 ¥ 10–3 m, Y = 7.9 ¥ 1010 N/m2 and D = 2650 kg/m3
From the relation, the fundamental frequency of quartz crystal
1 Y 1 7.9 ¥ 1010
f = =
2l D 2 ¥ 3 ¥ 10 -3 2650
= 9.1 ¥ 105 Hz
= 0.91 MHz
E xamplE 13 Calculate the natural frequency of iron of 0.03 m length, the density of iron is 7.23 ¥ 103 kg/m3
and Young’s modulus 116 ¥ 1010 N/m2.
Solution Given, l = 0.03 m, D = 7.23 ¥ 103 kg/m3 and Y = 116 ¥ 1010 N/m2
Formula used is
1 Y 1 116 ¥ 1010
f = =
2l D 2 ¥ 0.03 7.23 ¥ 103
f = 0.211 ¥ 106 Hz
= 0.21 MHz
E xamplE 14 An ultrasonic source of 0.67 MHz sends down a pulse towards sea bed which come back after
1 sec. Find out the depth of sea and the wavelength of pulse. The velocity of sound in sea water is 1690 m/sec.
Solution Given, f = 0.67 ¥ 106 Hz, t = 1 sec and v = 1690 m/sec
By using the formula
2h = vt and v = fl
where h is depth of the sea, we get
2 ¥ h = 1690 ¥ 1
h = 845 m
v 1690
and l= = = 0.00252 m
f 0.67 ¥ 106
E xamplE 15 Calculate the capacitance to produce ultrasonic waves of 106 Hz with an inductance of 1 Henry.
Solution Given, f = 106 Hz and L = 1 Henry
Formula used is
1
f =
2p LC
1
or C=
4p 2 Lf 2
1 10-12
C= 2 6 2
=
4 ¥ (3.14) ¥ 1 ¥ (10 ) 4 ¥ (3.14)2
= 0.0254 pF
E xamplE 16 A quartz crystal of thickness 1 mm is vibrating at resonance. Calculate the fundamental

frequency. Given Y for quartz = 7.9 ¥ 1010 N/m2 and r for quartz = 2650 kg/m3.
Solution The fundamental frequency of the vibration is given by
1 Y
f =
2l e
1 7.9 ¥ 1010
f=
2 ¥ 0.001 2650
= 2.72998 ¥ 106 Hz
The fundamental frequency of the quartz crystal = 2.730 ¥ 106 Hz = 2.73 MHz
o bJecTiVe T YPe Q uesTioNs
Q.1 If the period of a wave is decreased, then

(a) the amplitude of the wave decreases (b) the amplitude of the wave increases
(c) the frequency of the wave decreases (d) the frequency of the wave increases
Q.2 Which of the following frequency range is audible to the human ear?
(a) 50-100 Hz (b) 500-1000 Hz (c) 5000-10000 Hz (d) all of the above
Q.3 A property of sound which is most closely associated with the pitch of a musical note is
(a) amplitude (b) frequency (c) wave velocity (d) all of the above
Q.4 Two tones have the same amplitude. The statement which is true is
(a) the sounds must be equally loud to the ear
(b) the sounds must have the same pitch
(c) the sounds must have the same timbre
(d) none of the above must be true (though they may be true)
Q.5 When the amplitude of a sound pressure wave is doubled, the sound pressure level (in decibels)
(a) is doubled (b) is halved (c) decreases by 2 dB (d) increases by 6 dB
Q.6 ‘‘This note is higher in pitch than the note’’ is a statement about
(a) relative pitch (b) absolute pitch (c) binaural hearing (d) none of these
Q.7 Lying on the floor, you will exert pressure on the floor, compared to standing on the floor
(a) more (b) less (c) the same (d) not known
Q.8 A mass oscillates on a spring. As the mass passes through the point of equilibrium,
(a) kinetic energy is maximum and potential energy is maximum
(b) kinetic energy is minimum and potential energy is maximum
(c) kinetic energy is maximum and potential energy is minimum

(d) kinetic energy is minimum and potential energy is minimum
Q.9 The difference between the time each of your ears hears a sound can help you judge of the sound.
(a) direction (b) pitch (c) loudness (d) speed
Q.10 As the temperature of the air increases, the speed of sound in air
(a) increases (b) decreases (c) does not change (d) none of these
Q.11 Which of the following is the false statement?
(a) sound waves are longitudinal waves (b) sound can travel through a vacuum
(c) light travels much faster than sound (d) transverse waves in a guitar string are different from
sound waves
Q.12 The Doppler shift explains
(a) why a sound grows quieter as we move away from the source
(b) why the siren on a police car changes pitch as it races past us
(c) the phenomenon of beats
(d) sound diffraction
Q.13 The two tones of 440 Hz and 444 Hz are played. The beat frequency is
(a) 440Hz (b) 444 Hz (c) 442 Hz (d) 4 Hz
Q.14 Which of the following will not affect the fundamental frequency of vibration of a string?
(a) changing the amplitude of vibration (b) changing the tension of the string
(c) changing the length of the string (d) changing the density of the string
Q.15 Compared to the velocity of a 400 Hz sound through air, the velocity of a 200 Hz sound through air is
(a) twice as great (b) one-half as great (c) the same (d) times larger
Q.16 In general, sound travels fastest through
(a) gases (b) liquids (c) solids (d) vacuum
Q.17 In a sound wave, at a place where there is a node in the air pressure wavy, the air molecule displacement
(a) has a node (b) has an antinode
(c) oscillates between node and antinode (d) none of the above
Q.18 When a sound wave passes from air into water, it changes direction. This phenomenon is known as
(a) refraction (b) diffraction (c) reflection (d) polarisation
Q.19 A sound pressure level of 110 decibels would be considered
(a) very loud (b) average for speaking
(c) very soft (d) nothing can be said
Q.20 In damped harmonic motion, the following quantity decreases.
(a) amplitude (b) frequency (c) period (d) velocity
Q.21 If the displacement of vibrating particles is perpendicular to the direction of propagation of the wave,
the wave is said to be
(a) linear (b) transverse (c) longitudinal (d) standing
Q.22 There are two organ pipes of the same length. One has one end closed, while the other has both ends
open. The one with a closed end will emit sound of
(a) a higher frequency (b) a lower frequency
(c) the same frequency (d) double amplitude
Q.23 When a stretched wire (fixed at both ends) is vibrating in the second harmonic, there is/are
(a) one node (b) two nodes (c) three nodes (d) infinite nodes
Q.24 Two pure tones cause resonance in different positions along the basilar of membrane. These tones have
different
(a) amplitude (b) frequency (c) timbre (d) intensity
Q.25 Overtones have wavelengths, compared to the fundamental
(a) longer (b) shorter (c) the same (d) times larger
Q.26 In order to double the wavelength of a sound wave, you should only
(a) double its amplitude (b) double its frequency
(c) halve its amplitude (d) halve its frequency
Q.27 When two sine waves that are 180° out of phase are added together, the amplitude of the sum is
(a) always zero (b) always less than the amplitude of either wave
(c) equal to the amplitude of the smaller (d) always less than the amplitude of the larger wave
wave (e) always greater than the amplitude of the smaller wave
Q.28 A sound wave has sound intensity level SIL = 50 dB. Recall that SIL = 10 log ([I/10–12 W/m2)]). The
intensity I of this wave, in W/m2, is therefore
(a) 50 (b) 5 (c) 10–5 (d) 10–7 (e) 10–10
Q.29 A sound wave with SIL = 50 dB is reflected by a cloth-covered wall that absorbs 75% of its intensity.
The SIL of the reflected wave is
(a) 75 dB (b) 47 dB (c) 44 dB (d) 25 dB (e) 12.5 dB
Q.30 Light and sound are both waves; yet we can hear a car that is coming from behind the corner of a
building before we can see the car. This is because
(a) sound travels faster than light (b) sound lsound > llight, sound diffracts more than light
(c) sound is not reflected by buildings (d) sound and light interfere, with sound winning out
Q.31 A moving locomotive is sounding its horn as it crosses a highway. There are people in all directions
from the locomotive – in front, in back, to the right and left. Compared to the ‘‘true’’ pitch, as heard
by the engineer, the horn’s pitch heard by these people is
(a) higher (b) lower
(c) the same for all of the people (d) higher for some, true for others, and lower for yet
others of the people
Q.32 The frequency of the note B4 is close to 500 Hz. The period of this vibration is
(a) 500 sec (b) 1 sec (c) 0.2 sec (d) 2 msec
(e) none of these
Q.33 A sine wave and a square wave cannot have the same
(a) loudness (b) wavelength (c) frequency (d) tone quality (e) pitch
Q.34 An electric bell is operating in a vacuum. We cannot hear the sound of the bell because
(a) air is needed to conduct the electric current to the bell
(b) the bell’s metal cannot vibrate in vacuum
(c) there is no air to conduct the vibrations to our ears
(d) the vacuum jar absorbs the sound
(e) the noise of the pump is louder than the noise of the bell
Q.35 The wavelength of ‘‘shortwave’’ radio waves is smaller than that of standard broadcast (AM) radio
waves. They both propagate at the same speed. This allows you to conclude that, compared to AM
waves, the ‘‘shortwaves’’ have
(a) lower frequency (b) longer period (c) higher frequency (d) smaller amplitudes
Q.36 When a sound wave enters from air into a metal, in which the speed of sound is much larger than in
air, it does not change its
(a) wavelength (b) frequency (c) speed (d) all of these (a-c) change
Q.37 Sound moves at 345 m/sec towards a rock wall, reflects, and returns (as an echo). The roundtrip takes
2 sec. How far away is the wall?
(a) 70 m (b) 170 m (c) 340 m (d) 345 m (e) 350 m
Q.38 Two identical sound sources differ in distance from the listener by 1/2 wavelength. The result will be
(a) no sound at the listener
(b) constructive interference
(c) sound which is twice as loud as one source
(d) beats
Q.39 Which of the following crystals show piezoelectric effect?
(a) NaCl (b) Barium Titanate (c) Diamond (d) Quartz
Q.40 The frequency of vibration of the D.C. magnetized rod in the magnetostriction generator is
(a) Equal to the frequency of alternating current
(b) Twice the frequency of alternating current
(c) Half the frequency of alternating current
(d) None
T rue or F alse
state whether true or false

Q.1 A steady tone played on a violin is an almost perfect sine wave.
Q.2 Different vowel sounds differ mainly in the relative frequency and amplitude of the first two formats.
Q.3 The threshold of hearing is at 0 phons for all frequencies.
Q.4 A standing wave remains constant, without any change in time whatever.
Q.5 The precedence effect enables us to hear the fundamental frequency of a complex wave, even when
that frequency is absent in the Fourier spectrum.
Q.6 In the well-tempered scale, only the octaves are perfect intervals.
Q.7 Light is a longitudinal wave, whereas sound is transverse.
Q.8 A triangle wave contains higher-frequency Fourier components than a sine wave of the same periodicity.
Q.9 In a CD player the disk rotates at a constant linear velocity.
Q.10 In a dynamic loudspeaker the sound is produced by vibration of a permanent magnet.
Q.11 In an audio system, AM-FM tuner, tape recorder, and CD player each requires its own separate
amplifier and loudspeaker.
Q.12 To make the acoustics of an auditorium more live, the wall, ceiling, and floor surfaces should be made
as sound absorbent as possible.
Q.13 In white noise all frequencies are present, and all have the same intensity.
Q.14 Attack transients help determine the tone quality of a musical note.
Q.15 The sound quality of a violin is due only to the resonances of its strings; the violin body has no
resonances of its own.
Q.16 Percussion instruments have only a single resonance.

Q.17 For two electrical devices connected to the same voltage, the one with the smaller resistance draws the
smaller current.
Q.18 In order to avoid interference between different AM stations, each station uses a carrier wave that is
different from that of any of the other stations.
Q.19 Magnetostriction oscillator can generate ultrasonic waves of single frequency.
Q.20 We can generate ultrasonic waves of frequency f, 2f, …, where f is the fundamental frequency.
Q.21 The magnitude of the Piezoelectric effect of a crystalline material does not depend on direction.
Q.22 The Piezoelectric effect of a material is something to do with the crystal symmetry.
Q.23 The Piezoelectric effect of a material depends on its crystal structure.
P racTice P roblems
Q.1 What is piezoelectric effect? Describe the construction of a piezoelectric oscillator for the production
of ultrasonic waves.
Q.2 Give the theoretical treatment of Sabine’s law. Define the term ‘period of revereberation’.
Q.3 Sketch a graph of pressure vs time for two sound waves that differ only in pitch. Sketch a graph of
pressure vs time for two sound waves that differ only in timbre.
Q.4 A mass on a spring is found to oscillate naturally at a frequency of 0.5 Hz. This mass-spring system
is then driven by an oscillator. Describe what happens as the frequency of the oscillator is varied from
0.2 Hz to 0.8 Hz.
Q.5 Sketch the first two normal modes of sound pressure in a tube open at one end, closed at the other end.
If the fundamental mode has a frequency of 440 Hz, what is the frequency of the other mode? Is it
harmonic?
Q.6 Ram is in a fire truck rushing toward the scene of a fire. Shyam is standing at the scene of the fire.
There is no wind. Who hears a higher pitch for the fire truck’s siren? Explain why.
Q.7 Explain what is meant by a restoring force? Why is it necessary for vibrations to occur?
Q.8 What is the wavelength of a 440 Hz sound in air, if the speed of sound in air is 340 m/s? Would the
wavelength be longer or shorter if the sound were passing through water?
Q.9 A wave pulse travels down the length of a wave machine like the one in the front of the lecture hall.
The pulse reflects from the end. Describe the difference you would notice between a wave machine
with the end free to move and a wave machine with the end fixed.
Q.10 Define diffraction. How would you demonstrate diffraction of sound waves?
Q.11 One sound is made up of equal amplitudes of 110 Hz, 220 Hz, and 440 Hz pure tones. A second sound
is made up of equal amplitudes of 110 Hz, 330 Hz, and 550 Hz pure tones. In what way(s) are these
two sounds the same? In what way(s) are these two sounds different? What is the term given to this
combining of pure tones to get a complex tone?
Q.12 Two pure tones are played, one at a constant frequency of 550 Hz, the other has a variable frequency.
Describe all the phenomena you hear as the frequency of the second tone is varied gradually from
550 Hz to 1100 Hz.
Q.13 Why not ultrasonics be produced by passing high frequency alternating current through a loud speaker?
Dielectrics 10
Learning Objectives
LO 1 Understand the concept of dielectric LO 5 Discuss dielectric loss and the
constant dependent factors
LO 2 Know about types of dielectrics and LO 6 Describe Lorentz field and validity
polarisation of dielectrics of Clausius-Mosotti equation for
LO 3 Learn about types of polarisation non-polar dielectrics of cubic crystal
structure
LO 4 Explain Gauss’s law in the presence of
dielectrics
Introduction
A dielectric is an insulating material in which all the electrons are tightly bound to the nuclei of the
atoms and there are no free electrons available for the conduction of current. Therefore, the electrical
conductivity of a dielectric is very low. The conductivity of an ideal dielectric is zero. On the basis of band
theory, the forbidden gap (Eg) is very large in dielectrics. Materials such as glass, polymers, mica, oil and
paper are examples of dielectrics. They prevent flow of current through them. Therefore, they can be
used for insulating purposes.
10.1 DieLectric cOnstant LO1

It is found experimentally that the capacitance of a capacitor is increased if the space between its plates
is filled with a dielectric material. To understand this fact, Faraday took two identical capacitors, one was
evacuated and the other was filled with dielectric material, as shown in Fig. 10.1.
Then these two capacitors were charged with a battery of same potential difference. He found that the charge
on the capacitor filled with dielectric is larger than that of the other filled with air. If C0 be the capacitance in
I
Figure 10.1
vacuum and C the capacitance when the space is filled with a dielectric material, then the dielectric constant
of the material
C
K=
C0
Thus, the dielectric constant of a material is the ratio of the capacitance of a given capacitor completely filled
with that material to the capacitance of the same capacitor in vacuum. In other words, the ratio of permittivity
of medium to that of the vacuum is also known as dielectric constant, i.e.,
e
K= = er
e0
This is also known as relative permittivity (er). It is found to be independent of the shape and dimension of
the capacitor.
10.2 types Of DieLectrics LO2

A molecule is a neutral system in which the algebraic sum of all the charges is zero. Based on the dipole
moment, the molecules of dielectrics are termed as non-polar and polar molecules. Accordingly these
dielectrics are referred to as non-polar and polar dielectrics.
10.2.1 non-polar Dielectrics
A ‘non-polar’ molecule is the one in which the centre of gravity of the positive
(protons) and negative charges (electrons) coincide. So such molecule does
not have any permanent dipole moment, as shown in Fig. 10.2a. Few common
examples of non-polar molecules are oxygen (O2), nitrogen (N2) and hydrogen (a) (b)
(H2). As mentioned earlier, the dielectrics having non-polar molecules are Figure 10.2
known as non-polar dielectrics.
10.2.2 polar Dielectrics
A polar molecule is the one in which the centre of gravity of the positive charges is separated by finite
distance from that of the negative charges. Unbalanced electric charges, usually valence electrons, of such
molecules result in a dipole moment and orientation. Therefore, these molecules possess permanent electric
dipole (Fig. 10.2b). Few examples of polar molecules are N2O, H2O and HCl. The dielectrics having polar
molecules are known as polar dielectrics.
Dielectrics 315
10.3 pOLarisatiOn Of DieLectrics LO2

When an electric field is applied to a dielectric material; it exerts a force on each charged E0
particle and pushes the positive charge in its own direction while the negative charge is displaced
in opposite direction, as shown in Fig. 10.3. Consequently, the centres of positive and negative
charges of each atom are displaced from their equilibrium positions. Such a molecule (or atom) Figure 10.3
is then called as induced electric dipole and this process is known as dielectric polarisation.
We consider a parallel plate capacitor which has vacuum initially between its plates. When it is charged with
a battery, the electric field of strength E0 is set up between the plates of the capacitor (Fig. 10.4a). If s and –s
are the surface charge densities of the two plates of the capacitor, then the electric field developed between
the plates is given by
s
E0 = (i)
e0
If now a slab of dielectric material is placed between the two plates of the capacitor (Fig. 10.4b), then it
becomes electrically polarised. Hence, its molecules become electric dipole oriented in the direction of the
field. Because of this the centre of positive and negative charges gets displaced from each other. Therefore, in
the interior of the dielectric as marked by dotted lines these charges cancel. However, the polarisation charges
on the opposite faces of the dielectric slab are not cancelled. These charges produces their own electric field
Ep, which opposes the external applied field E0. Under this situation, the net electric field in the dielectric is
given by
  
E = E0 - E p (ii)
Figure 10. 4
10.3.1 polarisation Density
The induced dipole moment developed per unit volume in a dielectric slab on placing it inside an electric field
is known as polarisation density. It is denoted by a symbol P. If p is induced dipole moment of individual
atom and N is the number of atoms in a unit volume, then polarisation density is
 
P = Np (iii)

The induced dipole moment of an individual atom is found to be proportional to the applied electric field E
and is given by
 
P = ae 0 E (iv)
where a is the proportionality constant and is also known as atomic

polarisability. From Eqs. (iii) and (iv), we get
 
P = N ae 0 E (v)
Suppose S is the area of each plate of the capacitor and d is the

separation between them (Fig. 10.5). Then the volume of the dielectric
slab is Sd. Since –qi and +qi are the induced charges developed on the
two faces of the dielectric slab, the total dipole moment of the slab
will be equal to qid. From the definition of the polarisation density
Total dipole moment
P= Figure 10.5
Volume of slab
qi d qi
= = =sp
Sd S
\ P = sp (vi)
On placing the dielectric material between the two plates of the capacitor, the reduced value of the electric
field may be evaluated as follows
s -sp s sp
E= = - (vii)
e0 e0 e0
P
or E = E0 - (viii)
e0
From Eq. (vii)
e0E = s – sp = s – P [ P = sp]
or s = e0 E + P (ix)

The quantity (e0E + P) is of special significance and is known as the electric displacement vector D given by
  
D = e0 E + P (x)
10.3.2 Relation between Dielectric Constant and Electric Susceptibility

The polarisation density of a dielectric is proportional to the effective value of electric field E and is
given by
 
P = ce 0 E (xi)
where c is constant of proportionality and is known as susceptibility of dielectric material.
Dielectrics 317
By using Eq. (viii), we get
ce 0 E
E = E0 - = E0 - c E
e0
or E0 = E(1 + c) or E0/E = 1 + c
or K=1+c [ K= E0/E]
10.4 types Of pOLarisatiOn LO3

The important types of polarisation are categorised as
under.
10.4.1 electronic polarisation
Under the action of an external field, the electron clouds of
atoms are displaced with respect to heavy fixed nuclei to a
distance less than the dimensions of the atom (Fig. 10.6).
This is called electronic polarisation, which does not
depend on temperature. The electronic polarisation is
represented as below
 
Pe = N a e E (i) Figure 10.6
10.4.2 ionic polarisation
This type of polarisation occurs in ionic crystals, for example in sodium chloride crystal. In the presence of an
external electric field, the positive and negative ions are displaced in opposite directions until ionic bonding
forces stop the process (Fig. 10.7). This way, the dipoles get induced. The ionic polarisation does not depends
upon temperature.
E
Absence of field Presence of field

(a) (b)
Figure 10.7
10.4.3 Orientation polarisation
This types of polarisation is applicable in polar dielectrics. In the
absence of an external electric field, the permanent dipoles are
oriented randomly such that they cancel the effects of each other
(Fig. 10.8a). When the electric field is applied, these dipoles tend
to rotate and align in the direction of the applied filed (Fig. 10.8b).
This is known as orientation polarisation, which depends upon
temperature. Figure 10.8
In view of all these polarisations, the total polarisation is the sum of the electronic, ionic and orientation
polarisations. This is given by
   
P = Pe + Pi + Po
10.5 Gauss’s Law in DieLectrics LO4

Gauss’s law states that “the surface integral of electric field vector over a closed surface is equal to 1/e0 times
the net charge enclosed by the surface”, i.e.,
  1
Ú E ◊ dS = e 0 q
   
or Ú E ◊ dS = q or
e0  Ú e 0 E ◊ dS = q
Let us consider a parallel plate capacitor without
dielectric, as shown in Fig. 10.9a. Then Gauss’s law
is written as
 
Ú E ◊ dS = q
e0 
or e0E0S = q Figure 10.9
q
or E0 = (i)
e0 S
where E0 is the electric field between the plates.
In the presence of dielectric material between capacitor plates, it is clear from Fig. 10.9b that the total charge
enclosed by Gaussian surface is (q – q¢), where q¢ is the induced charge in the dielectric material due to
polarisation. Then Gauss’s law says
 
Ú E ◊ dS = (q - q ¢)
e0  (ii)
or e0ES = q – q¢
q q¢
or E= - (iii)
e0 S e0 S
As we know that the relative permittivity is
E0 E0
K= or E= (iv)
E K
With the help of Eq. (i), we get
E0 q
E= = (v)
K K e0 S
Dielectrics 319
Using Eq. (v) in Eq. (iii) we obtain

q q q¢
= -
K e0 S e0 S e0 S
q
or = q - q¢ (vi)
K
The induced charge is therefore
Ê 1ˆ
q ¢ = q Á1 - ˜
Ë K¯
From the above relation it is clear that the induced charge q¢ is less than the free charge q.
q
By substituting the value of (q - q ¢ ) = in Eq. (ii), we get for Gauss’s law
K
  q   q q  
Ú E ◊ dS = K or Ú E ◊ dS = e 0 K = e or Ú e E ◊ dS = q
e0 
 
i.e., Ú D ◊ dS = q
This is the Gauss’s law in the presence of a dielectric.
10.6 DieLectric LOss LO5

When a dielectric material is placed in an alternating electric field (Fig. 10.10a), a part of the energy is
wasted, which is known as dielectric loss. This is because of the fact that the reversing nature of the field
causes the direction of the dipoles to reverse. The dielectric loss depends on the frequency and the mechanism
by which the polarisation is produced in the material. An ideal dielectric does not absorb electrical energy.
However, in a real dielectric, there is always a loss of some electrical energy.
Figure 10.10
Consider a parallel plate capacitor of capacity C, whose plates having area S are separated by a distance d. The
space between the plates of the capacitor is filled with dielectric material having permittivity e. The sinusoidal
voltage V of angular frequency w is applied to the capacitor. Then the current through the capacitor is
Q V CV V
I= + = +
t R t R
Since w is inversely proportional to the time period, the current I may be written as
V
I = w CV +
R
or I = Ic + Id (i)
where Ic and Id are conduction and displacement currents, respectively. From the above relation, it is clear that
there are two kinds of currents that flow through the dielectric.
The Current Ic, Id and I are plotted in Fig. 10.10b, from which it is clear that the resultant current I = ( I c2 + I d2 )1/2
V
lags behind the displacement current by an angle d. In an ideal dielectric R = •, which means I c = becomes
R
zero. In this situation the resultant current would be
I = Id = wCV
e0 S e e S
Now C = free space and 0 r for the capacitor filled with dielectric material. Therefore,
d d
e 0 e r SwV
I = Id = (ii)
d
Ic
The angle d is known as loss angle, which can be calculated from Fig. 10.10b. In DOAD, tan d =
Id
so Ic = Id tan d (iii)
Eqs. (ii) and (iii) yield
(we 0 e r SV ) tan d
Ic =
d
Thus, the real power loss in the dielectric materials is
Pl = VI c
we 0 e r SV 2
= tan d
d
we e ( Sd )V 2
= 0 r 2 tan d
d
2
ÊV ˆ
= 2p f e 0 e rV Á ˜ tan d
Ëd¯
= (2pe 0 )e r f VE 2 tan d
Pl = 5.54 ¥ 10 -11 VE 2 e r f tan d
The above expression shows that the power loss depends on the volume V of the dielectric, its dielectric constant
er, frequency f of the alternating field together with its amplitude E.
Dielectrics 321
10.7 ClauSiuS-MoSotti Equation LO6

The interaction between the atoms of gases can be neglected as the atoms are separated by sufficiently large
distance. So atoms of gases feel the same field as applied to them. However, in solids and liquids, the atoms
are closely surrounded (on all the sides) by other atoms which may be polarised under the action of an
external field. Hence, the internal intensity of the electric field at a given point of the dielectric is generally
not equal to the intensity of the applied field. The internal field is actually the electric field acting at the
location of a given atom and is equal to the sum of the electric field created by neighbouring atoms and the
applied field. In case of crystal possessing cubic symmetry, the internal field is given by
  
Ei = E + P/3e 0
This field is called Lorentz fieldand also some time referred to as local field. In linear and isotropic dielectric
the molecular dipole moment p is directly proportional to the internal field Ei, that is
 
p = a Ei
  
p = a ( E + P/3e 0 ) (i)
As mentioned earlier, the proportionality constant a here is known
 as molecular polarisability. In terms of
number of molecules per unit volume N, the polarisation density P is given by
   
P = N p = N a ( E + P/3e 0 ) (ii)

The above equation for P gives
 Na  (iii)
P= E
1 - N a /3e 0
In terms of the electric susceptibility c, the above equation is written as
 
P = e0 c E (iv)
A comparison of Eqs. (iii) and (iv) yields
Na
= e0 c (v)
1 - N a /3e 0
Along with the use of K = 1 + c as e = e0 (1 + c) and K = e/e0, Eq. (v) gives
3e 0 ( K - 1)
a=
N ( K + 2)
K - 1 Na
or =
K + 2 3e 0
The above equation is known as Clausius-Mosotti equation. This equation is valid for non-polar dielectrics
having cubic crystal structure. The Clausius-Mosotti equation is also known as Lorentz-Lorentz equation in
view of its application in optics.
Physical significance of clausius
Mosoth equation
10.7.1 physical significance
Clausius–Mossotti equation or relation seems to hold best for gases but gives reasonably good results for
many liquids and solids too. This is clear from the relation that it connects the relative permittivity K of a
dielectric medium to the polarizability a of the atoms (or molecules) constituting the dielectric. Since the
relative permittivity is a bulk or macroscopic property and the polarisability is a microscopic property of mat-
ter, the Clausius-Mosotti equation bridges the gap between a directly observable macroscopic property with
a microscopic molecular property.
S UMMARY
✦ A dielectric is an insulating material in which all the electrons are tightly bound to the nuclei of the
atoms and there are no free electrons available for the conduction of current. Therefore, the electrical
conductivity of a dielectric is very low. The forbidden gap (Eg) is very large in dielectrics. Materials
such as glass, polymers, mica, oil and paper are a few examples of dielectrics.
✦ A non-polar molecule is the one in which the centre of gravity of the positive charge (protons) and
negative charge (electrons) coincide. So such molecules do not have any permanent dipole moment.
Nitrogen (N2) and hydrogen (H2) are the examples of non-polar molecules.
✦ A polar molecule is the one in which the centre of gravity of the positive charges is separated by finite
distance from that of the negative charges. Unbalanced electric charges, usually valence electrons,
of such molecules result in a dipole moment and orientation. Therefore, these molecules possess
permanent electric dipole. Examples of polar molecules are N2O, H2O and HCl.
✦ An external electric field, when applied to a dielectric material, exerts a force on each charged particle
and pushes the positive charge in its own direction while the negative charge is displaced in opposite
direction. Consequently, the centres of positive and negative charges of each atom are displaced from
their equilibrium positions. Such a molecule (or atom) is then called as induced electric dipole and this
process is known as dielectric polarisation.
✦ The induced dipole moment developed per  unit volume in a dielectric on placing it inside an electric
field is known as polarisation density P . If N be  the number
 of atoms in a unit volume and a the
atomic polarisability, then polarisation density is P = N ae 0 E .
✦ Electric susceptibility c and the dielectric constant K are related as K = 1 + c.
✦ Polarisation is of three types, namely electronic polarisation, ionic polarisation and orientation
polarisation.

✦ Gauss’s law states that the surface integral of electric field vector E over a closed surface is equal to
 q
1/e0 times the net charge enclosed by the surface, i.e., Ú E ◊ dS = e 0
 1 1
✦ The energy stored in an electrostatic field E is u = e 0 KE 2 which takes the form u = e 0 E 2 in the
free space. 2 2
✦ The internal intensity of the electric field at a given point of the dielectric is generally not equal to
the intensity of the applied field. The internal field is actually the electric field acting at the location
Dielectrics 323
of a given atom and is equal to the sum of the electric field created by neighbouring atoms and the
K - 1 Na
applied field. In case of crystal possessing cubic symmetry, a relation = exists between the
K + 2 3e 0
dielectric constant K, atomic polarisability a and permittivity e0. This equation is known as Clausius-
Mosotti equation, which is also known as Lorentz-Lorentz equation in view of its application in optics.
S OLVED E XAMPLES
E xamplE 1 Two parallel plates having equal and opposite charges are separated by a 2 cm thick slab that has
dielectric constant 3. If the electric filed inside is 106 V/m. Calculate the polarisation and displacement vector.

Solution Given E = 106 V/m = 106 N/C, K = 3, e 0 = 8.85 ¥ 10-12 C2 N -1m -2
  
Formula used is D = e0E + P
 
Also D = e 0 KE
or D = 8.85 ¥ 10-12 ¥ 3 ¥ 106

= 2.655 ¥ 10-5 C/m 2
  
P = D - e0E
= 2.655 ¥ 10-5 - 8.85 ¥ 10-12 ¥ 106
= 1.77 ¥ 10-5 C/m 2
E xamplE 2 Two parallel plates have equal and opposite charges. When the space between them is evacuated
the electric intensity is 3 ¥ 105 V/m and when the space is filled with dielectric the electric intensity is 1.0 ¥ 105
V/m. What is the included charge density on the surface of the dielectric?
 
Solution Given E0 = 3 ¥ 105 V/ m, and E = 1 ¥ 105 V/m
Formula used is

  P  
E = E0 - or P = s p = e 0 ( E0 - E ) [ P = s p ]
e0
sp = 8.85 ¥ 10–12 [3 – 1] ¥ 105
sp = 1.77 ¥ 10–6 C / m2
E xamplE 3 Two parallel plates of capacitor having equal and opposite charges are separated by 6.0 mm
thick dielectric material of dielectric constant 2.8. If the electric field strength inside be 105 V/m, determine
polarisation vector, displacement vector and energy density in the dielectric.
Solution Given E = 105 V/m = 105 N/C and K = 2.8
    1
P = e 0 ( K - 1) E , D = e 0 KE and energy density = K e 0 E 2
2
P = 8.85 ¥ 10 ¥ (2.8 – 1) ¥ 10 = 1.593 ¥ 10 C / m2
−12 5 −6
= 1.6 ¥ 10–6 C / m2
D = 8.85 ¥ 10–12 ¥ 2.8 ¥ 105 = 2.478 ¥ 10–6 C / m2
= 2.5 ¥ 10–6 C / m2
1
Energy density = ¥ 2.8 ¥ 8.85 ¥ 10-12 ¥ (105 ) 2
2
= 12.39 ¥ 10–2 J / m3
= 0.124 J / m3
E xamplE 4 An isotropic material of relative permittivity er is placed normal to a uniform external electric
field with an electric displacement vector of magnitude 5 ¥ 10–4 C / m2. If the volume of the slab is 0.5 m3 and
magnitude of polarisation is 4 ¥ 10–4 C / m2, find the value of er and total dipole moment of the slab.
 
Solution Given D = 5 ¥ 10-4 C/ m 2 , P = 4 ¥ 10-4 C/ m 2 and V = 0.5 m3 er = K = ?
Formula used is
  
 D = e0E + P
  
E = ( D - P )/e 0
(5 - 4) ¥ 10-4
=
8.85 ¥ 10-12
= 1.13 ¥ 107 V/m
D 5 ¥ 10-4
or K = er = = =5
e0E 10-4
Total dipole moment p
P= =
Volume V
p = PV = 4 ¥ 10–4 ¥ 0.5 = 2.0 ¥ 10–4 C-m
E xamplE 5 Dielectric constant of a gas at N.T.P is 1.00074. Calculate dipole moment of each atom of the
gas when it is held in an external field of 3 ¥ 104 V/m.

Solution Given E = 3 ¥ 104 V/m = 3 ¥ 104 N/C and K = e r = 1.00074
Formula used is K = 1 + c
or c = K – 1 = 1.00074 – 1 = 0.00074
and polarisation density is
P = c e0 E = 0.74 ¥ 10–3 ¥ 8.85 ¥ 10–12 ¥ 3 ¥ 104
= 19.647 ¥ 10–11 C / m
No. of atoms of gas per cubic metre (N)
6.06 ¥ 1023
= = 2.7 ¥ 1025
22.4 ¥ 10-3
P 19.647 ¥ 10-11
Induced dipole moment of each atom (p) = =
N 2.7 ¥ 1025
or p = 7.27 ¥ 10–36 C-m
E xamplE 6 Determine the electric susceptibility at 0°C for a gas whose dielectric constant at 0°C is 1.000041.
Solution Given K = 1.000041 and T = 0°C
Dielectrics 325
Formula used is
K=1+c
or c = K – 1 = 1.000041 – 1 = 0.41 ¥ 10–4
= 4.1 ¥ 10–5
O BJECTIVE T YPE Q UESTIONS
Q.1 Dimension of atomic polarisability in SI units is

(a) Cm–2 (b) CV–1 m2 (c) CVm (d) none of them
Q.2 A non-polar molecule is the one in which the centre of gravity of positive and negative charges
(a) coincides (b) gets separated by 10–8 m
(c) gets separated by 1Å (d) none of these
Q.3 The net charge inside a dielectric before and after polarisation remains
(a) negative (b) positive (c) same (d) none of these
Q.4 In vacuum, electric susceptibility is
(a) less than 1 (b) greater than 1 (c) small but –ve (d) zero
Q.5 The electric susceptibility and dielectric constant are related as
(a) c = K ~ 1 (b) c = 1 + K (c) c = e0K (d) c = 1 + e0K
Q.6 Dimension of displacement vector in SI unit is
(a) C-m (b) C-m–1 (c) C-m–2 (d) C-m2
  
Q.7 The relation between three electric vectors E , D and P is
          1  
(a) D = e 0 E + P (b) D = e 0 ( E + P) (c) D = E + e 0 P (d) D = E+P
e0
Q.8 Polarisation density is
(a) dipole moment per unit area (b) surface charge density
(c) dipole moment per unit volume (d) (b) and (c)
Q.9 Which of the following relation reperesent the Clausius-Mosotti equation
K - 1 Na K +1 K + 2 Na
(a) = (b) = Na (c) = (d) None of these
K + 2 3e 0 K -2 K - 1 3e 0
Q.10 Which type of polarisation is depend on temperature
(a) Orientational (b) Electronic (c) Ionic (d) none of these
Q.11 Gauss’s law in dielectrics is
  1    
(a) Ú E ◊ dS = e 0 q (b) Ú D ◊ dS = q (c) Ú E ◊ dS = e 0 q (d) both (a) and (b)
Q.12 Dimension of permittivity of free space in SI unit is
(a) Fm–1 (b) Cm–1 (c) Hm–1 (d) dimensionless
Q.13 Dimension of relative permittivity in SI unit is
(a) Fm–1 (b) Cm–1 (c) Hm–1 (d) dimensionless
Q.14 What changes in the capacitance of a capacitor occur if the dielectric material between the plates of
the capacitor is replaced by air or vacuum?
(a) increases (b) decreases (c) remains same (d) none of these
Q.15 Dielectrics are the substances which are

(a) conductors (b) insulators (c) semi-conductor (d) none of these
Q.16 Which of the statement is true?
(a) Ionic polarisation decreases with increase of temperature
(b) Ionic polarisation does not depend on temperature
(c) Ionic polarisation increases with rise in temperature
(d) none of these
Q.17 Which of the following statement is true?
(a) Orientational polarisation decreases with rise in temperature
(b) Orientational polarisation increases with rise in temperature
(c) Orientational polarisation does not depends on temperature
(d) none of these
S HORT -A NSWER Q UESTIONS
Q.1 What are polar and non-polar molecules?

Q.2 Discuss different polarisation mechanisms in dielectrics?
Q.3 What is a dielectric?
Q.4 How long does polarisation of non-polar molecules last?
Q.5 State and prove Gauss law in dielectrics.
Q.6 What are dielectric losses?
Q.7 Discuss the behaviour of a dielectric in a.c. field.
Q.8 Write the Clausius-Mossotti equation.
Q.9 Will an atom having spherically symmetric charge distribution be polar or non-polar? Explain.
Q.10 What do you understand by polarisation of dielectric and dielectric susceptibility? Find the relation
between the two.
Q.11 Write note on
(i) Dielectrics (ii) Three electric vectors (iii) Dielectric losses.
P RACTICE P ROBLEMS
General questions
Q.1 What is a dielectric substance? Give examples. Discuss the importance of dielectrics.
Q.2 What are polar and non-polar molecules? Discuss the effect of electric field on polar dielectrics. What
is meant by polarisation of dielectric?
Q.3 Discuss different types of polarisations in dielectrics.
Q.4 What happens when a non-polar molecule is placed in an electric field? Define atomic dipole moment
and atomic polarisability. What are their dimensions? Give their S.I. units.
Q.5 What is atomic polarisability?
 Find a relation between dipole moment and atomic polarisability or

show that p = a E0 .
Dielectrics 327

 P
Q.6 Show that the electric field inside a polarised dielectric due to induced polarisation charge is E = -
where P is the polarisation density vector. e0
Q.7 Explain the terms dielectric polarisation, susceptibility, permittivity and dielectric coefficient. Derive
their inter-relation equation.
  
Q.8 Define and explain the three electric vectors P, E and D . Why electric field inside a dielectric
  
decreases due to polarisation? Show that D = e 0 E + P . Also give their units.
  
Q.9 Show that D = e 0 E + P , where the symbols have their usual meanings.
Q.10 What are three electric vectors in dielectrics? Name and find relation between them.
Q.11 What do you understand by polarisation of dielectric and dielectric susceptibility? Find the relation
between the two.
Q.12 Explain the phenomenon of polarisation of dielectric medium and show that K = 1 + ce. Here the
symbols have their usual meanings.
Q.13 Define the terms dielectric constant K and electric susceptibility ce. Prove the relation K = 1 + ce.
Q.14 Find the relation between induced charge and free charge when a dielectric material of dielectric
constant K is placed between the plates of a parallel plate capacitor.
-1
È sp ˘
Q.15 Prove that induced charge varies with the dielectric as K = Í1 - ˙ , where sp and sfree are the
Î s free ˚
induced and free surface charge densities, respectively. Hence show that for a metal K = •.

Q.16 The electric field between the plates of a parallel plate capacitor is E0 without dielectric. But if
dielectric of relative permittivity er is introduced between the plates what will the electric field be?
Q.17 What is the effect of temperature on the dielectric constant of a substance containing molecules of
permanent dipole moment?
Q.18 Derive a relation between electric susceptibility and atomic polarisability on the basis of microscopic
description of matters at atomic level.
Q.19 Derive Clausius-Mosotti relation for non-polar dielectrics.
Q.20 Discuss the effect of introducing a dielectric between the plates of a capacitor. Show that the capacitance
of a charged capacitor when a dielectric material of dielectric constant K is introduced between the
plates is given by K e 0 A .
d
Q.21 Explain why the introduction of a dielectric slab between the plates of a capacitor changes its
capacitance?
Q.22 State and prove Gauss’s law in dielectrics.
Q.23 Derive an expression for Gauss’s law in the presence
  of dielectric. Prove that divergence of displacement
vector is equal to density of free charge or — ◊ D = rfree . Also discuss integral form of Gauss’s law.
 
Q.24 Using Gauss’s law in dielectric medium, show that — ◊ D = rfree , where symbols have their usual
meanings.
Q.25 Explain the mechanism contributing to dielectric polarisation. Discuss the behaviour of a dielectric in
an alternating field.
1  
Q.26 Show that the electrostatic energy per unit volume in a dielectric is D ◊ E , where symbols have their
usual meanings. 2
Q.27 Deduce an expression for energy stored in dielectric in electrostatic field.
Electromagnetism 11
Learning Objectives
LO 1 Understand charge density, del LO 9 Discuss scalar and vector potential,
operator, gradient, divergence and curl continuity equation, maxwell’s
LO 2 Explain fundamental theorem of equation in differential and integral
calculus and for gradient form
LO 3 Know about Gauss’s theorem and its LO 10 Evaluate significance of Maxwell’s

functional correlations equation, Maxwell’s displacement
current and correction in Ampere’s
LO 4 Discuss Stokes theorem and its use in Law
calculating electric field and electric
potential LO 11 Learn about electromagnetic wave
propagation, Transverse nature
LO 5 Know how to derive Poisson’s and EMW, Maxwell’s equation in Isotropic
Laplace’s equations dielectric medium/conducting
LO 6 Know about capacitor and their medium, electromagnetic energy
important configurations density
LO 7 Understand the concepts of magnetic LO 12 Elaborate on Poynting vector and
flux density, magnetic field strength and Poynting theorem, waveguide, coaxial
Ampere’s circuital law cables
LO 8 Explain electrostatic boundary
conditions
Introduction
Ordinary matter is made up of atoms which have positively charged nuclei and negatively charged electrons
surrounding them. Charge is quantized in terms of the electronic charge –e. One electronic charge e is
equal to 1.602 ¥ 10–19 Coulombs. One Coulomb of charge is the charge which would flow through a
220 W light bulb (220 Vac) in one second. Do you know two charges of one Coulomb each separated
by one meter would repel each other with a force of about a million tons! The separation of charges
produces electric field, whereas the motion of charges generates current and hence the magnetic field.
When these fields are time varying they are coupled with each other through the Maxwell’s equations.
Electromagnetism 329
With the help of the Maxwell’s equations, we can derive wave equation, based on which the propagation
of electromagnetic waves can be investigated in different media.
11.1 Charge Density LO1

If a charge is distributed continuously in a medium, it can be expressed in terms of a physical quantity known
as charge density. There are three types of charge densities, namely linear charge density, surface charge
density and volume charge density.
11.1.1 Linear Charge Density (l)

If a charge is distributed continuously on a linear conductor, the charge on its unit length is called the linear
charge density. It is generally represented by l and is measured in the units of C / m. If l be the total length of
a conductor and l the linear charge density, then the total charge on the conductor will be
l
q = Ú l dl
0
For a uniform charge distribution, l is constant, and is given as l = q / l. Here q is the total charge, given by
q = ll.
11.1.2 surface Charge Density (s)

If the charge is distributed over a surface, the charge on the unit area of the conductor is called the surface
charge density (s) and its unit is C / m2.
If the surface charge density at a point of the conductor is s, then the charge contained in a small element of
area dS will be s dS. Therefore, the total charge on the surface of the conductor,
q = Ú s dS
s
If the charge distribution is uniform, then the value of s will be constant, and is given by
s = q/S
where S is total area of the surface and q is the total charge, given by q = sS.
11.1.3 Volume Charge Density (r)

If the charge is distributed in the volume of a conductor, the charge contained in a unit volume of the conductor
is called the volume charge density (r) and its unit is C / m3.
If the volume charge density at a point of the volume of a conductor is r, then the charge contained in a small
element of the volume dV will be r dV. Therefore, the total charge contained in the conductor,
q = Ú r dV
v
For the uniform distribution of the charge, r will be constant, and is given by
r = q/V
or q = rV
4
where q is the total charge and V is the total volume, which is equal to p r 3 for a spherical shape of the
conductor of radius r. 3
11.2 DeL OperatOr LO1


The del operator is the differential operator, which is represented by — and is given by — = iˆ ∂ + ˆj ∂ + kˆ ∂ .

∂x ∂y ∂z
It is not a vector in itself, but when it operates on a scalar function it provides the resultant as a vector. For

example, when — is operated on a scalar function F(x, y, z), we get
 ∂F ˆ ∂F ˆ ∂F
—F = iˆ + j +k
∂x ∂y ∂z
 
—F does
 not mean a multiplication of — with F rather it is an instruction to differentiate. Here we should say
that — is a vector operator that acts upon F.

 del operator — can act in different ways. For example, when it acts on a scalarfunction F, the resultant
The
— F is called  the
 gradient of a scalar function F. When it acts
 on a vector function A via the dot product, the
resultant is — ◊ A, which is called
  the divergence of a vector A . When
 it acts on a vector function A via the cross
product, the resultant is — ¥ A, which is called the curl of a vector A. Finally, the Laplacian of a scalar function
∂2 ∂2 ∂2
F is written as —2F together with —2 = + + .
∂x 2 ∂y 2 ∂z 2
The del operator can be written in spherical polar coordinate system (Fig. 11.1a) as
 ∂ 1 ∂ 1 ∂
— = rˆ + qˆ + fˆ
∂r r ∂q r sin q ∂f
In cylindrical coordinate system (Fig. 11.1b), it is given by

 ∂ 1 ∂ ∂
— = sˆ + fˆ + zˆ
∂s s ∂f ∂z
z
z
P(s, f, z)
P(r, q, f)
z
q
r
y
y
f f
x x
(a) (b)
Figure 11.1
11.3 graDient LO1

If we think of the derivative of a function of one variable we notice that it simply tells us how fast the function
varies if we move a small distance. It means the gradient is the rate of change of a quantity with distance. For
example, temperature gradient in a metal bar is the rate of change of temperature along the bar. However, for
a function of three variables the situation is more complicated, as it will depend on what direction we choose
to move. For a function F(x, y, z) of three variables, we obtain from a theorem on partial derivative
∂F ∂F ∂F
dF = dx + dy + dz
∂x ∂y ∂z
Here dF is a measure of changes in F that occurs when we alter all three variables by small amounts dx, dy and dz.
The above expression for dF can be written in terms of a dot product of vectors as
 
dF = —F ◊ d l
 ∂F ˆ ∂F ˆ ∂F
where —F = iˆ + j +k is nothing but gradient of F. Clearly the gradient is a vector quantity, i.e.,
∂x ∂y ∂z
it has both magnitude and direction. The meaning of gradient becomes clearer when we write
 
dF = —F ◊ d l
 
or dF = |—F ||d l | cos a
  
where a is the angle between —F and d l . Now we fix the distance dl, i.e., magnitude |d l |, and look around in
various directions for the maximum change in F. Clearly the maximum change in F takes place in the direction
a = 0. It means dF is largest when we move in the direction of —F. This can also be said that the gradient —F points
in the direction of maximum increase of the function F. Hence, the gradient of a scalar field F is a vector quantity
which represents both the magnitude and the direction of the maximum space rate of increase of F.
The gradient of F in Cylindrical coordinate system is written as
 ∂F 1 ∂F ˆ ∂F
—F = sˆ + f+ zˆ
∂s s ∂f ∂z
In spherical polar coordinate system, the gradient of F is written as
 ∂F 1 ∂F ˆ 1 ∂F ˆ
—F = rˆ + q+ f
∂r r ∂q r sin q ∂f
11.4 DiVergenCe LO1

  
As mentioned earlier, the divergence of a vector field A is represented by — ◊ A. So it is given by
  ∂A ∂Ay ∂Az
—◊ A = x + +
∂x ∂y ∂z
Clearly the divergence
 of a vector field is a scalar. Also, the divergence of a scalar cannot be obtained, as the
dot product of — with a scalar is not possible.
  
In order to make clear the meaning of the divergence,
 we consider the net flux  Ú A ◊ dS of a vector field A
from a closed surface S. The divergence of A is defined as the net outward flux per unit volume over a closed
 
surface. Actually the divergence of A at a given point is a measure of how much the vector A spreads out, i.e.,
diverges, from that point. Fig. 11.2a shows that the divergence of a vector field at point O is positive, as the vector
spreads out. However, in Fig. 11.2b, the vector converges and hence the divergence at O is negative. In Fig. 11.2c,
we can notice that the divergence of a vector field is zero, as the magnitude of vectors remains the same.
O O
(a) (b)
(c) (d)
Figure 11.2

By now the difference between vector field and vector would have been clear to you. We call A as vector field
because its values at different points are different, which are vectors with different magnitudes. For example,
in Fig. 11.2d the magnitudes of the vectors get increased as we move towards right. Hence, the divergence of
such a vector field is not zero but positive.

The divergence of a vector A in cylindrical coordinate system can be written as
  1 ∂( sAs ) 1 ∂Af ∂Az
—◊ A = + +
s ∂s s ∂f ∂z
In spherical polar coordinate system, it is written as
  1 ∂ 2 1 ∂ 1 ∂Af
—◊ A = 2 (r Ar ) + (sin q Aq ) +
r ∂r r sin q ∂q r sin q ∂f
  
Finally, we mention that the vector field A is said to be solenoidal or divergenceless if — ◊ A = 0.
11.5 CurL LO1

  
As mentioned earlier, the curl of vector field A is represented by — ¥ A. So it is given by
  Ê ∂A ∂Ay ˆ Ê ∂A ∂A ˆ Ê ∂Ay ∂Ax ˆ
— ¥ A = iˆ Á z - ˜ + ˆj Á x - z ˜ + kˆ Á -
Ë ∂y ∂z ¯ Ë ∂z ∂x ¯ Ë ∂x ∂y ˜¯
Clearly the curl is a vector quantity, i.e., it has both

 direction and magnitude. Also, it is evident that we cannot
have the curl of a scalar as the cross product of — with a scalar is meaningless.
In order to make clear the meaning of curl, we consider the circulation of a vector X
   
field A around a closed path, i.e.,  Ú A ◊ dl . It is evident that the curl of A is a
rotational vector. Its magnitude would be the maximum circulation of A per unit
area. Its direction is the normal direction of the area when the area is oriented so as Z
to make the circulation maximum.
 Actually curl of A at some point O is a measure
of how much the  vector A curls around the point O. So, you may notice that the curl
Y
of vector field A in Fig. 11.2 is zero. However, in Fig. 11.3, the curl is finite and is
Figure 11.3
pointing in the x-direction as per right hand rule.

The curl of a vector A in cylindrical coordinate system can be written as
  Ê 1 ∂Az ∂Af ˆ Ê ∂A ∂A ˆ 1Ê ∂ ∂A ˆ
—¥ A=Á - ˜ sˆ + Á s - z ˜ fˆ + Á ( sAf ) - s ˜ zˆ
Ë s ∂f ∂z ¯ Ë ∂z ∂s ¯ s Ë ∂s ∂j ¯
In spherical polar coordinate system, it is written as
  1 È ∂ ∂Aq ˘ 1 Ê 1 ∂Ar ∂ ˆ ˆ 1È ∂ ∂Ar ˘ ˆ
—¥ A= Í ∂q (sin q Af ) - ∂f ˙ rˆ + r ËÁ sin q ∂f - ∂r (rAf )˜¯ q + r ÍÎ ∂r (rAq ) - ∂q ˙˚ f
r sin q
Î ˚
  
Finally, we mention that the vector field A is said to be irrotational or potential if — ¥ A = 0.
11.6 FunDamentaL theOrem OF CaLCuLus LO2
If g(x) be a function of one variable, then dg would represent infinitesimal change in g (x) when we move
Ê dg ˆ
from x to x + dx. This change in g(x) is given by dg = Á ˜ dx. If we move from point a1 to point a2, then the
Ë dx ¯
total change in g(x) can be obtained by using the fundamental theorem of calculus. This theorem states that
a2
Ê dg ˆ
Ú ÁË dx ˜¯ dx = g (a2 ) - g (a1 )
a1
It means the total change in the function can be obtained by simply subtracting the values of the function at
the points a2 and a1. For the movement from point a1 to point a2, these points can be treated as the end points.
Ê dg ˆ
In view of this, the fundamental theorem says that the integral of a derivative Á ˜ over an interval a1 Æ a2
Ë dx ¯
is given by the value of the function at the end points. The end points represent the boundaries.
Since we have three types of derivatives, namely gradient, divergence and curl in vector calculus, there are three
fundamental theorems for these derivatives. The fundamental theorem for divergence is also known as Gauss’s
theorem or Green’s theorem. Similarly, the fundamental theorem for curl is also known as Stokes’ theorem.
11.7 FunDamentaL theOrem FOr graDient LO2

We know that the gradient is defined only for a scalar function. So we consider a scalar  function
 of three
variables F(x, y, z). As discussed earlier, the change in this function is defined as dF = —F ◊ d l . If we move
from point a1 to point a2 then the total change in F can be calculated using fundamental theorem of calculus.
Therefore, the total change in F in moving from a1 to a2 is given by
a2
 
Ú —F ◊ dl = F (a2 ) - F (a1 )
a1
This is called the fundamental theorem for gradient, according to which the integral of a derivative over end
points a1 and a2 is given by the value of the function at the boundaries, i.e., points a1 and a2. Here it can be
noticed that the integral is line integral, the derivative is the gradient and the boundaries are the points a1 and a2.
Moreover, it can be seen that the integral or the total change in function F is independent of path taken from
a1 to a2. Also, if we take both the end points same, i.e., we evaluate the close integral, then the total change
in the function F comes out to be zero (as the beginning and end points are identical).
11.8 gauss’s Or green’s theOrem LO3

According to this theorem,
   
Ú (—◊ F )dV = Ú F ◊ dS
v s

Here dV is the volume element and dS is the surface element. In the same manner as we stated fundamental
theorem for gradient, the Green’s theorem states that the integral of a derivative (here the divergence) over
a region (here the volume) is equal to the value of the function at the boundary (here the surface). Since the
boundary of a volume is always a closed surface, the R.H.S. is the integral over closed surface.
Evidently this theorem converts the volume integral into the surface integral. Therefore, this theorem is very
useful in the situations where it is difficult to calculate the volume integral.
1
The net outward electric flux through any closed surface drawn in an electric field is equal to times the
total charge enclosed within the surface. It is expressed as e0
  1 1
Ú E ◊ dS = e 0
f= Âq = e Q
s 0
where Q is the sum of all the charges present within the surface. The charge outside of the surface is not
counted, as the lines entering and leaving the surface due to this charge are the same in number. Therefore the
flux f due to the charge q sitting outside the surface is expressed as
 
Ú E ◊ dS = 0
f=
s
11.8.1 Differential Form of gauss’s theorem

When a charge is distributed over a volume such that r is the volume charge density, then the charge enclosed
by the surface enclosing the volume is given by
q = Ú r dV (i)
v
  q
Substituting this expression of q in Ú E ◊ dS = e 0 , we get
s
 
Ú E ◊ dS = Ú rdV
e0  (ii)
s v
According to Gauss’s divergence theorem, we can convert the surface integral into the volume integral as
  
Ú E ◊ dS = Ú div E dV (iii)
s v
Therefore, from Eqs (ii) and (iii), we have


e 0 Ú div EdV = Ú r dV
s v
Since the above equality is true for every volume, the integrands of left and right sides should be equal, i.e.,

e0divE = r
 
— ◊ (e0E ) = r (iv)
 
—◊ D = r (v)
  
where, D is the electric flux density, given by D = e0E . Eq. (v) is the differential form of Gauss’s theorem.
11.9 stOkes’ theOrem LO4

Stokes’ theorem states that the integral of the curl of a vector function over a patch of surface is equal to the
value of the function at the perimeter of the patch. So here the derivative is the curl, region is the surface and
the boundary is the perimeter of the patch of the surface. Therefore,
    
Ú (— ¥ F ) ◊ dS = Ú F ◊ dl
s c
Clearly, this theorem converts the surface integral into the line integral. Here the L.H.S. is the surface integral
whereas the R.H.S. is the closed line integral. So a point of confusion is that which way we should go around,
i.e., clockwise or anticlockwise when  we integrate the line integral. Moreover,
 we should know about the
direction of the surface element dS . For example, for a closed surface dS points outwards normal but for
an open surface which way is out? In order to overcome this confusion, we apply right  hand rule. So if our
fingers point in the direction of line integral, then the thumb gives the direction of dS .
  
Based on the statement of Stokes’ theorem, we can make some more observations. For example, Ú (— ¥ F ) ◊ dS
s
does not depend on the shape of the surface rather it depends on the boundary line. Also for a closed surface,
  
Ú (— ¥ F ) ◊ dS = 0 as the boundary line shrinks down to a point.
s
11.10 eLeCtriC FieLD anD eLeCtriC pOtentiaL LO4

Electric
 field is the region around a charge in which another charge  experiences a force. The electrostatic
field E is a special kind of field whose curl is always zero, i.e., — ¥ E = 0. Since any vector whose curl is
zero is equal to the gradient of some scalar quantity, we make use of this property to introduce the scalar
    
quantity as the electric potential V. In order to find this, we use the Stokes’ theorem Ú (— ¥ E ) ◊ dS =  Ú E ◊ dl .
This gives
 
Ú E ◊ dl = 0 (i)
It means the line integral of E  from point a1 to point a2 will be the same for all the paths between these points.
Hence, the line integral of E is independent of path. Since changing the path would not alter the value of
integral, we can define a function, say V, such that
r
 
V = - Ú E ◊ dl (ii)
a
The differential form of the above equation is written as
 
E = –—V (iii)
Here V is called the electric potential. Actually all the potentials are relative and there is no absolute zero
potential. However, convention is that the potential is zero at infinite distance from the charge. In view of this,
the lower limit a in Eq. (ii) is called as a standard reference point where V is zero. The upper limit is nothing

but the point where V is to be calculated. So V depends only on the point r .
11.10.1 superposition principle

According to the original principle of superposition of electrodynamics, the total force F on a charge q (test
charge) is equal to the vector sum of the forces due to all the source charges (considering them individually).
It means    
F = F 1 + F 2 + F 3 + ... (iv)
  
Since F = qE , from Eq.(iv) we find the following for the electric field E
   
E = E 1 + E 2 + E 3 + ... (v)
If we write a for the common reference point, the above equation can be written as
r r
    r   r  
- Ú E ◊ dl = - Ú E1 ◊ dl - Ú E2 ◊ dl - Ú E3 ◊ dl - º
a a a a
(vi)
   
or —V = —V1 + —V2 + —V3 + ... (vii)
Now it is clear from Eq. (vii) that
V = V1 + V2 + V3 + ... (viii)

The above equation reveals that the potential V at a given point r is the sum of the potentials due to all the
charges. It means the electric potential also satisfies the principle of superposition
 and the sum is simply an
ordinary sum. However, from Eq. (v) it is clear that in case of the electric field E this sum is the vector sum.
11.11 pOissOn’s anD LapLaCe’s equatiOns LO5

For deriving these equations, we start with the following Gauss’s law for a linear medium
 
—◊ D = r
 
or — ◊ eE = r (i)

Here e is basically the free charge density (volume) and D is the electric displacement (talked in detail later).
 
Since E = –—V, the above equation for a homogeneous medium (where e is constant) can be written as
—2V = –r/e (ii)
This equation is called as Poisson’s equation. For a charge free region, i.e., where r = 0, the Poisson’s
equation takes the form
—2V = 0. (iii)
This equation is known as Laplace’s equation. This equation is much useful in solving electrostatic problems
where a set of conductors are maintained at different potentials; for example, capacitors and vacuum tube
diodes.
Using the expressions for Laplacian operator —2 in Cartesian, cylindrical and spherical coordinate systems,
we can write Laplace’s Eq. (iii) in these coordinate systems, respectively, as
∂ 2V ∂ 2V ∂ 2V
+ + =0
∂x 2 ∂y 2 ∂z 2
1 ∂ Ê ∂V ˆ 1 ∂ 2V ∂ 2V
Ás ˜+ + =0
s ∂s Ë ∂s ¯ s 2 ∂f 2 ∂z 2
1 ∂ Ê 2 ∂V ˆ 1 ∂ Ê ∂V ˆ 1 ∂ 2V
Á r ˜ + Á sin q ˜ + =0
r 2 ∂r Ë ∂r ¯ r 2 sin q ∂q Ë ∂q ¯ r 2 sin 2 q ∂f 2
11.12 CapaCitOr LO6

In order to construct a capacitor, we need at least two conductors that carry equal but opposite charges. If it is
so, then flux lines leaving one conductor terminate at the surface of the other conductor. These conductors are
also known as the plates of the capacitor, which may be separated by a dielectric or simply by free space. The
ratio of the magnitude of the charge on one of the plates (say q) and the potential difference (say V) between
the plates is known as the capacitance of the capacitor (represented by C). So
q
C=
V
Some important capacitor configurations with two conductors are parallel plate capacitor, coaxial capacitor
and spherical capacitor. These are discussed below.
11.12.1 parallel plate Capacitor

This type of capacitor has two plane conductors, which are parallel to each other. If the area of each plate be
S and they are separated by a distance d, then the capacitance of this capacitor is given as
S
C=e
d
Here e is the permittivity of the dielectric, which is filled between the plates. If a potential difference V is
applied between the plates then the energy stored in the capacitor would be
1 q2
U= CV 2 =
2 2C
11.12.2 Coaxial Capacitor
Coaxial capacitor is simply a coaxial cable. This is also referred to as coaxial cylindrical capacitor. If L
be the length of the coaxial conductors, and the radii of inner and outer conductors be rin and rout, then the
capacitance of the capacitor is obtained as
2pe L
C=
r
ln out
rin
Here also e is the permittivity of the dielectric filling the space between the two conductors.
11.12.3 spherical Capacitor
As name suggests, in this case the two conductors are in the form of spheres and these are concentric. Let the
radius of inner sphere be rin and of the outer sphere be rout. Also these spheres are separated by a dielectric
medium of permittivity e. Then the capacitance of this type of capacitor is obtained as
4pe
C=
1 1
-
rin rout
Æ
11.13 Magnetic Flux Density (B ) LO7
When a magnetic material is placed in an external magnetic field, it gets magnetised. The magnetism thus
produced in the material is known as induced magnetism and this phenomenon is referred to as magnetic
induction. The magnetic lines of force inside such magnetised materials are called magnetic lines of induction.
The number of  magnetic lines of induction crossing unit area at right angles to the flux is called the magnetic
flux density B. Its unit is the Tesla which is equal to 1 Wb / m2.
Æ
11.14 Magnetic FielD strength (H ) LO7
As mentioned earlier, a magnetic material becomes magnetised when placed in a magnetic field. The actual
magnetic field inside the material is the sum of external field and the field due to its magnetisation.

 B    
H= - M or B = m0 ( H + M )
m0
Magnetic field strength at a point in a magnetic field is the magnitude of the force experienced by a unit
pole situated at that point. The SI unit, corresponding to force of 1 Newton, is the A / m. The CGS unit,
corresponding to a force of 1 dyne is the Oersted which is equal to 79.6 A / m.
11.15 ampere’s CirCuitaL Law LO7

Ampere’s circuital law in magnetostatics  is analogous to the Gauss’s law used in electrostatics. This law says
that the line integral of magnetic field B around any closed loop is equal to m0 times the net current I flowing
through the area enclosed by the loop i.e.,
 
Ú B ◊ dl = m0 I
Here, m0 is the permeability of the free space.
Proof: Consider a long straight conductor carrying a current I. By Biot-Savart law, the magnitude of the
magnetic field at a point O, at a distance r from the conductor, is given by
m 2I
B= 0 (i)
4p r 
Let us draw a circle with a radius r taking C as centre around the current carrying conductor Fig. 11.4. B will
be the same in magnitude at all points on  this circle.
 Again we consider a circle element of length dl at the
point O. From the figure it is clear that d l and B are in the same direction.
 
\
Ú B ◊ dl = Ú Bdl cos q
→
B →
dl
C O
Ú dl
= B [ q = 0∞] r
m0 2 I I
= 2p r = m0 I
4p r
 
Ú B ◊ dl = m0 I Figure 11.4

This is the required Ampere’s circuital law. This law can be written in terms of volume current density J , if
    
we apply Stokes’ theorem 
 
Ú B ◊ dl = Ú (— ¥ B) ◊ dS . 
 
Since I = Ú J ◊ dS , we get — ¥ B = m0 J from  Ú B ◊ dl = m0I. This is the another form of Ampere’s law.
11.16 eLeCtrOstatiC BOunDary COnDitiOns LO8
Consider the situation where the electric field exists in a region, which has two different media with
permittivities as e1 (in region 1) and e2 (in region 2). Then the conditions, which should be satisfied by the
field at the interface separating the two media or at the common boundary of these media, are called boundary
conditions. It is seen that when we cross a boundary surface charge s, the electric field does not remain
continuous and it always undergoes a discontinuity.
E2
Amperian
Loop Gaussian
µ2 E2N Pillbox
2 e2
E2T A B
s
∆s Ds
E1T D C
1 Dw e1
µ1 E1 E1N
Figure 11.5

We
 can calculate the amount by which  the electric field E changes at such a boundary, shown in Fig. 11.5.
E 1 is the field in the region 1 and E 2 is the field in the region 2. These
 fields can be decomposed into two
components, out  of which one is tangential to the boundary (say E T) and the other is perpendicular to the
boundary (say E N). So we can write
     
E 1 = E 1N + E 1T and E 2 = E 2N + E 2T
Now we can apply Ampere’s and Gauss’s laws for calculating the amount of discontinuity. For example,
 
Ú E ◊ dl = 0 for the closed path ABCDA, whose length is Dw and the width is Ds, follows
Ampere’s law 
Ds Ds Ds Ds
0 = E2T Dw - E2 N - E1N - E1T Dw + E1N + E2 N (i)
2 2 2 2
In the limit Ds Æ 0, i.e., when the width of the loop is small so that we are well close to the boundary, the
above equation says
E2T Dw = E1TDw
or E2T = E1T (ii)
  
Since D = eE , the above equation reads for the displacement vector D T
D2T D1T
= (iii)
e2 e1
Eq. (ii) says that the tangential component of the electric field remains continuous across the boundary,
 as its
values just below and just above the boundary are equal. It means the electric field component
 E T undergoes
no change on the boundary. However, you  can see from Eq. (iii) that the field component DT undergoes some
change across the boundary. It means DT is discontinuous across the boundary.
 
In order to check the continuity of the normal component E N of the field E , we select Gaussian pillbox of
  
area S (upper and lower surfaces) and the thickness Ds. Now we apply Gauss’s law  Ú D ◊ dS = q and obtain
the following under the limit of Ds Æ 0
D2N S – D1N S = sS
or D2N – D1N = s (iv)
Here
 s is the free charge density placed at the boundary. It is clear from Eq. (iv) that the normal component of
D is discontinuous and this discontinuity
 amounts to the free charge density s. If there is no free charge, the
normal component of the field D will be continuous at the boundary, as in the absence of s Eq. (iv) follows
D2N – D1N = 0 (v)

From the above equation, we can find the condition for the electric field component E N as
e2E2N – e1E1N = 0 (vi)
 
The above equation shows that the normal component E N of the field E will be discontinuous at the boundary.
If we write Eqs (ii), (iii), (iv) and (vi) together, these equations represent the boundary conditions, named
dielectric – dielectric boundary conditions.
The boundary conditions are useful in finding the electric field on one side of the boundary if the field on
the other side is given.  In addition
 to this, we can determine the refraction of the  electric
 field across the
boundary. If the field E 1 (or D1) in the region 1 makes an angle a1 and the field E 2 (or D2) in the region 2
makes an angle a2 with the normal to the boundary, then from Eq. (ii) we get
E1 sin a1 = E2 sin a2 (vii)
Similarly, Eq. (vi) yields
e1E1 cos a1 = e2E2 cos a2 (viii)
From Eqs (vii) and (viii), we finally get

tan a1 tan a 2
=
e1 e2
tan a1 tan a 2
or = [ e1 = e0er1, e2 = e0er2] (ix)
e r1 e r2
Here er1 is the dielectric constant of medium 1 and er2 is the dielectric constant of medium 2. The above
equation is known as law of refraction of the electric field at the interface which is free of charge. This is also
clear from Eq. (ix) that in general a boundary between two dielectrics produces bending of the flux lines. This
is attributed to unequal polarisation charges accumulated on the sides of the boundary.
Now we can discuss some special cases of the above boundary conditions. For example, if the dielectric in
the medium 2 is replaced with a conductor whose  conductivity
 is infinite. In this case, it is obtained that the
field does not exist within the conductor, i.e., E = 0. Since E = –—V, a potential difference between any two
points in the conductor cannot exist. In other words, we say that a conductor is an equipotential body. Under
this situation, only the normal component of the electric field survives, i.e., DT = e0erET = 0, DN = e0erEN = s.
Similarly, when the dielectric in medium 1 is replaced with a conductor and the dielectric in medium 2 is
replaced with free space (er = 1), then the boundary conditions take the form DT = e0ET = 0, DN = e0EN = s.
11.17 sCaLar anD VeCtOr pOtentiaLs LO9

  
As mentioned
  of electrostatic field E , i.e., — ¥  E = 0, introduces a scalar potential V such
 earlier, the zero curl
that E = –—V. When we analyze — ◊ B = 0, we find that the field B can be written as a curl of another vector
quantity (say A), i.e.,
  
B = —¥ A (i)

Since a field can be completely determined if we know its divergence as well as its curl, the divergence of A
remains to be explored. In this context, with the use of Eq. (i), Ampere’s law reads
         
— ¥ B = — ¥ (— ¥ A) = —(— ◊ A) – —2A = m0J (ii)
 
It is clear that Eq.(ii) will resemble Poisson’s equation, if — ◊ A = 0. This condition is known as Coulomb
gauge. With the application of this condition, the Ampere’s law simply yields
 
—2A = –m0J (iii)

The solution of the above equation can be obtained if the current density J vanishes at infinity. Then the
solution comes out to be

  m0 J
4p Ú r
A(r ) = dX (iv)

Here dX is the volume element and the vector  A is called magnetic vector potential. Like the electric scalar
potential V, the magnetic vector potential A cannot be uniquely  defined as we can add to it another vector
whose curl is zero. This addition does not change the field B. On  the  other side, it is a point of observation
that we cannot introduce a magnetic scalar potential U such that B = –—U. The reason is that it is incompatible
with Ampere’s law, since the curl of a gradient is always zero.
11.18 COntinuity equatiOn LO9

The continuity equation says that the total current flowing out of some volume must be equal to the rate of
decrease of the charge within that volume, if the charge is neither being created nor lost. Since the charge
is flowing, we consider that the charge density r is a function of time. The transportation of the charge
constitutes the current, i.e.,
dq d
dt dt Úv
I= = r dV (i)
Here, it is assumed that the current is extended in space of volume V closed by a surface S. The net amount of 
charge which crosses a unit area normal to the directed surface in unit time is defined as the current density J .
This current density J is related to the total current I flowing through the surface S as
 
I= Ú J ◊ dS (ii)
s
Here the integral is over closed surface, as the surface bounding the volume is closed surface. From Eqs (i)
and (ii), we have
  dq - d
Ú J ◊ dS = - dt = dt Ú rdV (iii)
s v
The minus sign above is needed in view of decreasing charge r in the volume V. So
  ∂r
Ú J ◊ dS = - Ú ∂t dV (iv)
s v
From Gauss’s divergence theorem, we have
  
Ú J ◊ dS = Ú (div J )dV
s v
 ∂r
or Ú (div J ) dV = - Ú ∂t dV (v)
v v
Since the Eq. (v) holds good for any arbitrary volume, we can put the integrands to be equal. Then
 ∂r
div J + =0 (vi)
∂t
This is the continuity equation.
In case of stationary currents, i.e., when the charge density at any point within the region remains constant,
but the charges are moving.
∂r (vii)
= 0,
∂t
  
so that div J = 0 or —◊ J = 0

which expresses the fact that there is no net outward flux of current density J . Here the situation is the same
as shown in Fig. 11.2c for zero divergence.
11.19 maxweLL’s equatiOns: DiFFerentiaL FOrm LO10
When the charges are in motion, the electric and magnetic fields are associated with this motion which will
have variations in both the space and the time. These electric and magnetic fields are inter related. This
phenomenon is called electromagnetism which is summarised by the set of equations, known as Maxwell’s
equations. The Maxwell’s equations are nothing but are the representation of the basic laws of electromagnetism.
 
First Maxwell’s equation is the Gauss’s law of electrostatics, i.e., — ◊ E = r/e 0 . Second Maxwell’s
 
equation is the Gauss’s law of magnetostatics, i.e., — ◊ B = 0 . Faraday’s law of electromagnetic induction,

  -∂B
i.e., — ¥ E = , is called Maxwell’s third equation. Fourth Maxwell’s equation is the modified Ampere’s
∂t 
   ∂E
circuital law, i.e., — ¥ B = m0 J + m0e 0 .
∂t
Now we look for the convenient form of Maxwell’s equations while we are working with materials (having
permittivity e and permeability m) that are subject to electric and magnetic polarizations. Electric polarization

P provides bound charges with volume density eb, given by
 
rb = - — ◊ P (i)

Similarly, a magnetic polarization or magnetization M results in a volume bound current density, given by
  
Jb = — ¥ M (ii)
 
Equations (i) and (ii) represent the static case of uniforms polarization P and uniform magnetization M .

However, any change in polarization P gives rise to the polarization current density, given by

 ∂P
JP = (iii)
 ∂t 
Since J P satisfies the continuity equation, it is evident that J P is essential to account for the conservation of

bound charge. On the other hand, a changing magnetization M does not lead to any analogous accumulation
of charge or current. We do not have direct control on the bound charge and current. Hence, it is advisable
to reformulate Maxwell’s equations such that these make explicit  reference only to those sources which we
control directly. These are the free charges (rf) and currents (J f ). The total volume charge density can be
written as
r = rf + rb
 
= rf – — ◊ P , (iv)
whereas the total volume current density is written as

       ∂P
J = J f + Jb + J P = J f + — ¥ M + (v)
∂t
In view of total charge density r, Gauss’s law reads
  1  
— ◊ E = ( r f - — ◊ P)
e0
  
— ◊ (e 0 E + P) = r f
  
Recall that e 0 E + P = D . Hence, the first Maxwell’s equation in materials takes the form
 
— ◊ D = rf (vi)
Similarly, the modified Ampere’s circuital law yields

 
  Ê   ∂P ˆ ∂E
— ¥ B = m0 Á J f + — ¥ M + ˜ + m 0e 0
Ë ∂t ¯ ∂t

 Ê B ˆ  ∂  
—¥Á - M ˜ = J f + (e 0 E + P )
Ë m0 ¯ ∂t
 
   ∂D Ê  B ˆ
—¥ H = Jf + ÁË H = m - M ˜¯ (vii)
∂t 0
The four Maxwell’s equations in terms of free charges (density rf) and currents are written as
 
— ◊ D = rf First equation
 
—◊B=0 Second equation

  ∂B
—¥E=- Third equation
∂t

   ∂D
—¥ H = Jf + Fourth equation
∂t
 
However, for the sake of simplicity we shall write r for rf and J for J f in the derivation of Maxwell’s
equations, unless specified.
In free space, dielectric medium or conducting medium, the first and fourth Maxwell’s equations assume
different forms. For example, in free space and dielectric medium, free charge rf = 0 and free current density

J f = 0. Hence, the
 
First equation in free space yields — ◊ D = 0
  
fi — ◊ (e 0 E ) = 0 ( P = 0 in free space)
 
or — ◊E=0
However, the first equation in dielectric medium gives
     
— ◊ D = 0 fi e— ◊ E = 0 ( D = e E )
 
or —◊E =0
In conducting medium, any free charge resides on its surface, i.e., rf = 0 in the medium. Hence, the first equa-
 
tion again gives — ◊ E = 0 . This can also be understand as follows.
The continuity equations for free charges reads
∂r f  
+—◊Jf =0
∂t
∂r f  
or = -— ◊ J f (viii)
∂t
Since we want to see what happens when a free charge is given to a conductor, we find Eq. (viii) in terms of
   
rf by using Ohm’s law J f = s E (s is conductivity) and Gauss’s law — ◊ D = r f . Hence
∂r f  
= - — ◊ (s E )
∂t
∂r f s  
fi =- — ◊ (e E )
∂t e
∂r f s
fi =- rf (ix)
∂t e
1 ∂r f s
Eq. (ix) is written as = - and is integrated to get
r f ∂t e
-s t
r f (t ) = r f (0)e e (x)
Here rf(0) is the initial charge given to the conductor. For good conductors s ª •; means rf Æ 0 very quickly.
This proves that the charge will flow out to the edges of conductor within very less time. This characteristic
time is given by
e
t= (xi)
s  
So this is clear that the Maxwell’s first equation reads — ◊ E = 0 in free space, dielectric and conducting
medium.

   ∂D
Now we discuss different forms of Maxwell’s fourth equation — ¥ H = J f + . In free space
    ∂t
J f = 0 and D = e 0 E ( P = 0), Hence,

  ∂E
— ¥ H = e0
∂t

  ∂E
or m 0 — ¥ H = m 0e 0
∂t

 ∂E
or — ¥ B = m 0e 0 (xii)
∂t
  
In dielectric medium J f = 0 and D = e E . Hence

  ∂E
—¥H =e
∂t

  ∂E
or — ¥ ( m H ) = me
∂t

  ∂E
or — ¥ B = me (xiii)
∂t
   
In conducting medium, J f = s E and D = e E .
Hence,

   ∂E
— ¥ H =sE + e (xiv)
∂t
In view of the above discussion, this is clear that second and third Maxwell’s equations remain unchanged in
all types of the media.
11.19.1 Derivation of maxwell’s First equation


Let us consider a surface S bounding
 a volume V in a dielectric medium, which is kept in the E field. The
application of external field E polarises the dielectric medium and charges are induced, called bound charges
or charges due to polarisation. The total charge density at a point in a small volume element dV would then be
(r + rp), where rp is the polarisation charge density (the same as rb̂), given by rP = –div P and r is the free
charge density at that point in the small volume element dV.

Thus, the total charge density at that point will be r – (divP ). Then Gauss’s theorem can be expressed as
   1 
Ú E ◊ dS = Ú (div E )dV = e 0 Ú (r - div P) dV
s v v
 
or e 0 Ú (div E )dV = Ú ( r - div P ) dV
v v
 
or Ú div(e 0 E + P) dV = Ú rdV
v v
  Æ
The quantity (e0E + P ) is denoted by a quantity D , called the electric displacement. Therefore,

Ú (div D)dV = Ú rdV
v v
Since this equation is true for all the arbitrary volumes, the integrands in this equation must be equal, i.e.,

divD = r
 
or —◊ D = r
This is the Maxwell’s first equation.
   
When
 the medium is isotropic, the three vectors D, E and P are in the same direction and for small field E ,
D is proportional to E , i.e.,
 
D = eE
where e is called the permittivity of the dielectric medium. The ratio e / e0 is called the dielectric constant of
the medium.
11.19.2 Derivation of maxwell’s second equation

Since the magnetic lines of force are either closed or go off to infinity, the number of magnetic lines of force

entering any arbitrary surface is exactly the same as leaving it. It means the flux of magnetic induction B
across any closed surface is always zero, i.e.,
 
Ú B ◊ dS = 0
s
Transforming the surface integral to volume integral using Gauss’s divergence theorem, we have
  
Ú B ◊ dS = Ú (div B)dV = 0
s v
The integrand in the above equation should vanish for the surface boundary as the volume is arbitrary.
Therefore   
div B = 0 or — ◊ B = 0
This is the Maxwell’s second equation.
11.19.3 Derivation of maxwell’s third equation

According to Faraday’s law, the emf induced in a closed loop is given by

∂f ∂B  ∂  
Eemf = - = -Ú ◊ dS = - Ú B ◊ dS
∂t s
∂t ∂t s
 
Here the flux f = Ú B ◊ dS where S is any surface having the loop as boundary. The emf (Eemf) can also
s
be found by calculating the work done in carrying a unit charge completely around the loop. Thus,
 
Eemf = Ú E ◊ dl
c

Here E is the intensity of the field associated with the induced emf. On equating the above two equations, we get

  ∂B 
Ú E ◊ dl = - Ú ∂t ◊ dS
c s
According to Stokes’ theorem, the line integral can be transformed into surface integral with the help of
    
Ú E ◊ dl = Ú (— ¥ E ) ◊ dS . Therefore
c s 
   ∂B 
Ú (— ¥ E ) ◊ dS = - Ú ∂t ◊ dS
s s
This equation must be true for any surface whether small or large in the field. So the two vectors in the
integrands must be equal at every point, i.e.,

  ∂B
—¥E=-
∂t
This is the Maxwell’s third equation.
11.19.4 Derivation of maxwell’s Fourth equation

According to the Ampere’s law, the work done in carrying a unit magnetic pole once around a closed arbitrary
path linked with the current is expressed by
 
Ú H ◊ dl = I
c
   
or
Ú H ◊ dl = Ú J ◊ dS
c c
As per Stokes’ theorem,
    
Ú H ◊ dl = Ú (— ¥ H ) ◊ dS
c s
Therefore,
    
Ú ¥ H ) ◊ dS = Ú J ◊ dS
( —
s s
This gives
  
—¥ H = J
The above relation is derived on the basis of Ampere’s law, which holds good only for the steady current.
However, for the changing electric fields, the current density should be modified. The difficulty with the
above equation is that, if we take divergence of this equation, then
  
div(— ¥ H) = div J

fi 0 = div J [Since divergence of a curl = 0]

fi div J = 0
which conflicts with the continuity equation, as
 ∂r
div J = -
∂t
Therefore, Maxwell
 realised that the definition of the current density is incomplete and suggested to add
another density J ¢. Therefore
  
curl H = J + J ¢
Now, taking divergence of the above equation, we get
  
div (curl H) = div J + div J ¢
 
or 0 = div J + div J ¢
  ∂r
div J ¢ = - div J =
∂t
using continuity equation
Since,  
r = —◊ D
 ∂  
div J ¢ = (—◊ D)
∂t

   ∂D
—◊ J ¢ = —◊
 ∂t
 ∂D
Hence J ¢ =
∂t
Therefore, the Maxwell’s fourth equation can be written as

   ∂D
—¥H = J +
∂t
The last term of R.H.S. of this equation is called Maxwell’s correction and is known as displacement current
density. The above equation is called modified Ampere’s law for unsteady or changing current which is
responsible for the electromagnetic fields.
11.20 maxweLL’s equatiOns: integraL FOrm LO10

There are situations where the integral form of Maxwell’s equations is useful. Therefore, now we derive these
equations in integral form.
11.20.1 Maxwell’s First equation

Differential form of the Maxwell’s first equation is
 
—◊ D = r (i)
On integrating Eq. (i) over a volume V, we have
 
Ú (—◊ D)dV = Ú rdV
v v
Using Gauss’s divergence theorem, the above equation reads
 
Ú D ◊ dS = Ú rdV = q
s v
 
or Ú D ◊ dS = q
s
Here q is the net charge contained in the volume V and S is the surface bounding the volume V. This integral
form of the Maxwell’s first equation says that the total electric displacement through the surface S enclosing
a volume V is equal to the total charge contained within this volume.
This statement
 can also be put in the
 following form: The total outward flux corresponding to the displacement 
vector D through a closed surface S is equal to the total charge q within the volume V enclosed by the surface S .
11.20.2 Maxwell’s second equation

Differential form of the Maxwell’s second equation is
 
—◊ B = 0 (ii)
Exactly in a manner adopted above, we can show that
 
Ú B ◊ dS = 0
s
 
which signifies that the total outward flux of magnetic induction B through any closed surface S is equal to zero.
11.20.3 Maxwell’s third equation

Differential form of the Maxwell’s third equation is

  ∂B (iii)
—¥E=-
∂t

On integrating Eq. (iii) over a surface S bounded by a closed path, we have

   ∂B 
Ú ( — ¥ E ) ◊ dS = Ú ∂t ◊ dS
s s
Converting surface integral into line integral by Stokes’ theorem, we get

  ∂  
Ú E ◊ dl = - ∂t Ú B ◊ dS
c s
which signifies that the electromotive force around a closed path is equal to the time derivative of the magnetic
flux through any closed surface bounded by that path.
11.20.4 Maxwell’s Fourth equation

Differential form of the Maxwell’s fourth equation is

   ∂D
—¥H = J +
∂t (iv)
Exactly, in a manner adopted above, we can have this equation in the following form

  Ê  ∂D ˆ 
Ú H ◊ dl = Ú ÁË J + ∂t ˜¯ ◊ dS
c s
The above equation signifies that the magnetomotive force around a closed path is a measure of the conduction
current plus the time derivative of the electric flux through any surface bounded by that path.
11.21 signiFiCanCe OF maxweLL’s equatiOns LO10

Maxwell’s equations represent concisely the fundamentals of electricity and magnetism. From them one can
develop most of the working relationships in the field.
Maxwell’s first equation represents the Gauss’s law for electricity which says that the electric flux out of
any closed surface is proportional to the total charge enclosed within the surface. The integral form of this
equation finds applications in calculating electric fields around charge objects. It is consistent with Coulomb’s
law when applied to the electric field of a point charge. The area integral of the electric field gives a measure
of the net charge enclosed. However, the divergence of the electric field gives a measure of the density of
sources.
As mentioned, the area integral of a vector field determines the net source of the field (function). The integral
 
Ú B ◊ dS = 0 of the Maxwell’s second equation says that the net magnetic flux out of any closed surface
form 
s
is zero. This is because the magnetic flux directed inward toward the south pole, of a magnetic dipole kept in
any closed surface, will be equal to the flux outward the north pole. Therefore, the net flux is zero for dipole
 
sources. If we imagine a magnetic monopole source, the area integral  Ú B ◊ dS would have some finite value.
s
Because of this and since the divergence of a vector field is proportional to the density of point source, this
form of the Gauss’s law for magnetic field simply says that there are no magnetic monopoles.
The Maxwell’s third equation when written in the integral form states that the line integral of the electric
field around a closed loop is equal to the negative of the rate of change of the magnetic flux through the area
enclosed by the loop. The line integral basically is the generated voltage or emf in the loop. Therefore, the
physical interpretation of Maxwell’s third equation is that the changing magnetic field induces electric field.

For static electric field E , the second term of the R.H.S. of the Maxwell’s fourth equation vanishes and
then the integral form of this equation says that the line integral of the magnetic field around a closed loop
is proportional to the electric current flowing through the loop. This form of the Maxwell’s equation is
useful for calculating the magnetic field for simple geometries. However, this equation more specifically
reveals that the changing electric field induces magnetic field. This seems complimentary to the meaning
of the Maxwell’s third equation. Therefore, they together yield the formulation of electromagnetic fields or
electromagnetic waves, where both electric and magnetic fields propagate together and the change in one
field induces the other field.
11.22 maxweLL’s DispLaCement Current anD

LO10
COrreCtiOn in ampere’s Law
When a current flows in a conductor, magnetic  field is produced around it. A relation between the conduction
current
 (I) and the magnetic field vector (B ) was given by Ampere, according to which the line integral of the
field B along a closed curve in the magnetic field is equal to m0 times the current I flowing in the conductor.
   
Ú B ◊ dl = m0 I = m0 Ú J ◊ dS (i)
c s
   
or Ú (curl B) ◊ dS = m0 Ú J ◊ dS (Using Stokes’ theorem)
s s
  
or Ú (curl B - m0 J ) ◊ dS = 0
s
For arbitrary surfaces, the above integral is true, so we get
 
curl B – m0J = 0
 
or curl B = m0J (ii)
 
or div(curl B) = div(m0J )

or 0 = m0div(J ) [ divergence of curl of a vector field is always zero]

or div J = 0 (iii)
 ∂r  ∂r  ∂r
Using continuity equation div J + = 0 or div J = - , we get div J = 0 only if = 0 , i.e., for static
∂t ∂t ∂t
charge density.
Thus, the Ampere’s law given by Eq. (i) is valid only for the steady state conditions, i.e., for the fields that
do not vary with time. This law is not valid for the time varying fields, such as charging and discharging of a
∂r
capacitor, where π0.
∂t
Ampere’s law was modified by Maxwell for the time varying fields. His concept was based on the Faraday’s
law of electromagnetic induction, according to which a changing magnetic field produces an electric field.
On the basis of the fact that the magnetic field around a conductor is produced by the current flowing in it,
Maxwell hypothesized that changing electric field should also induce a magnetic field. A changing electric
field is equivalent to a current called displacement current (Id) which flows as long as the electric field is
changing. The displacement current produces the magnetic field the same way as the conduction current (I).
 
Thus, the total
 magnetic field (B) will be the sum of the two terms (a) due to current I(B1) and (b) due to
current Id(B2), i.e.,
  
B = B1 + B2
  
\ curl B = curl B1 + curl B2
  
or curl B = m0J + m0J d
  
or curl B = m0(J + J d) (iv)

where J d is the displacement current density. Also, in the integral form
 
Ú B ◊ dl = m0 I + m0 I d
c
 
or Ú B ◊ dl = m0 ( I + I d ) (v)
c
  dfB df
In analogy to the Faraday’s law of induction Ú E ◊ dl =-
dt
, I d should correspond to e 0 E . With this
dt
Eq. (v) can be written as c
  Ê df ˆ
Ú B ◊ dl
= m0 Á I + e 0 E ˜
Ë dt ¯
(vi)
c
 
dfE dE dD 
Thus, Id = e0 = e0 S =S = SJ d (vii)
 dt dt dt
where D is the electric displacement vector and S is the area.
Value of Jd can also be determined by taking divergence of Eq. (iv), i.e.,
   
div curl B = div (m0J + m0J d) = 0 [ div curl B = 0]
  ∂r È  ∂r ˘
or div J d = - div J = ÍÎ div J + ∂t = 0˙˚
∂t

 ∂  Ê ∂D ˆ 
div J d = (div D) = div Á
or
∂t Ë ∂t ˜¯ [ div D = r]

\
 ∂ D (viii)
Jd =
∂t
Therefore, the modified form of Ampere’s law is

 Ê  ∂D ˆ (ix)
curl B = m0 Á J + ˜
Ë ∂t ¯

 Ê  ∂D ˆ  
or curl H = Á J + ˜ ( B = m0H)
Ë ∂t ¯

  ∂E
or curl H = J + e 0 (x)
∂t
The Eqs (v), (ix) or (x) represent the Maxwell’s fourth equation in different form which is nothing but the
modified form of Ampere’s law.
11.23 eletroMagnetic (eM) Wave ProPagation in Free sPace LO11

Maxwell’s equations for free space are given as follows

div E = 0 (i)

div H = 0 (ii)

 ∂H (iii)
curl E = - m0
∂t

 ∂E
curl H = e 0 (iv)
∂t

   ∂       ∂ È ∂E ˘
Taking curl of Eq. (iii), we get — ¥ (— ¥ E ) = - m0 (— ¥ H ) or —(—◊ E ) - —2 E = - m0 Íe 0 ˙
∂t ∂t Î ∂t ˚
   
Here we have used Eq. (iv) for the value of — ¥ H. Now from Eq. (i) — ◊ E = 0. Hence

 ∂2 E
—2 E - m0e 0 2 = 0
∂t

This is the wave equation governing the field E . In view of the dimensions of (m0e0)–1/2 as of velocity (say,
v), we can write this as

 1 ∂2 E
—2 E - 2 2 = 0 (v)
v ∂t

Similarly the curl of Eq. (iv) gives rise to the wave equation for the field H as

2
 ∂2 H
— H - m 0e 0 2 = 0
∂t

 1 ∂2 H
or —2 H - 2 2 = 0 (vi)
v ∂t
The plane wave solutions of Eqs. (v) and (vi) may be written as
        
E (r , t ) = E0 ei ( k ◊ r - w t ) and H (r , t ) = H 0 ei ( k ◊ r - w t )
  
where w is the angular frequency of the variation of the fields E and H and k is the wave vector which tells
the direction of propagation of the fields or wave. The ratio w/k gives the phase velocity of the wave.
 Ê ∂ ∂ ˆ ∂ ˆˆ 
Now, — = Á iˆ + j + k ˜ , E0 = ( E0 xiˆ + E0 y ˆj + E0 z kˆ)
Ë ∂x ∂y ∂z ¯
 
k = (k xiˆ + k y ˆj + k z kˆ), r = ( xiˆ + yjˆ + zkˆ)
 
fi k ◊ r = (kxx + kyy + kzz)
   
\ curl E = — ¥ [ E0ei ( k ◊ r - wt ) ]
  
= — ¥ {( E0eik ◊ r )e - iwt }
 i(k x + k y + k z )
= — ¥ {[( E0 xiˆ + E0 y ˆj + E0 z kˆ)e x y z ]e - iwt }
Here note that i appeared in exponential term is such that i = -1, whereas iˆ is the unit vector along the
x-direction. When we solve the above equation with the help of expansion of curl, we obtain
 i ( k x + k y + k z z ) - iw t
curl E = i{iˆ[ E0 z k y - E0 y k z ] + ˆj[ E0 x k z - E0 z k x ] + kˆ[ E0 y k x - E0 x k y ]}e x y e
  i ( kˆ ◊ r - w t )
= i[k ¥ E0 ]e
  
or curl E = i[k ¥ E ] (vii)
Here we have used
   
E (r , t ) = E0 ei ( k ◊ r - w t )
Using Eqs. (iii) and (vii), we obtain

∂H  
- m0 = i[k ¥ E ]
∂t
L.H.S. of the above equation is written as
∂  i ( k ◊ r - w t ) 
- m0 [H 0e ] = iwm0 H
∂t
  
Hence k ¥ E = wm0H (viii)
Similarly, it can be shown using Eq. (iv) that
  
k ¥ H = –we0E (ix)

From
 Eq. (viii) it is obvious that  the magnetic field vector H
 is perpendicular to both the
 propagation
 vector
k and the electric field vector E and according to Eq. (ix) E is perpendicular to both k and H. Therefore, it
may be concluded that the electric  and magnetic vectors are normal to each
 other as well as to the direction
of propagation of the wave or E , H and direction of wave propagation k form a set of orthogonal vectors.
Further, we can prove that  the electromagnetic field or wave travels at the speed of light c in free space. For
this, the cross product of k with Eq. (viii) gives
    
k ¥ (k ¥ E ) = wm0(k ¥ H)
      
k (k ◊ E ) – k2E = wm0[–we0E ] [Putting the value of k ¥ H from Eq. (ix)]
   
Since k and E are perpendicular to each other, k ◊ E = 0 and the above equation reads
 
k2E – w2m0e0E = 0

(k2 – w2m0e0)E = 0
This relation between w and k is known as dispersion relation.

Since E cannot be zero for the wave, k2 – w2m0e0 = 0
w 1
fi = = 3 ¥ 108 m/sec = c, the speed of light.
k m0 e 0
w
Therefore, the phase velocity of the electromagnetic wave is equal to the speed of light c in free space or
vacuum. k
11.24 transVerse nature OF eLeCtrOmagnetiC waVes LO11


Transverse
 nature of electromagnetic waves can be proved with the help of Maxwell’s equations  div E = 
0
and div H = 0 for free space. Using different relations as discussed in the previous section div E and div H
can be calculated as follows.
  Ê ∂ ∂ ˆ ∂ ˆˆ  i ( k ◊ r - w t )
—◊ E = Á iˆ + j + k ˜ ◊ E0 e
Ë ∂x ∂y ∂z ¯
Ê ∂ ∂ ˆ ∂ ˆˆ i(k x + k y + k z - wt )
= Á iˆ + j + k ˜ ◊[( E0 x iˆ + E0 y ˆj + E0 z kˆ)e x y z ]
Ë ∂x ∂y ∂z ¯
∂ i(k x + k y + k z - wt ) ∂ i(k x + k y + k z -wt ) ∂ i(k x + k y + k z - wt )
= [ E0 x e x y z ] + [ E0 y e x y z ] + [ E0 z e x y z ]
∂x ∂y ∂z
i(kx x + k y y + kz z - wt )
= ( E0 x ik x + E0 y ik y + E0 z ik z )e
 
= i (k x E0 x + k y E0 y + k z E0 z )ei ( k ◊ r - w t )
   
= i (k ◊ E0 )ei ( k ◊ r - w t )
     
= ik ◊ E (  E0ei ( k ◊ r - wt ) = E )
 
\ For free space, — ◊ E = 0
   
ik ◊ E = 0 or k ◊ E = 0
 
It means the wave vector k is perpendicular to E .
   
Similarly, from — ◊ H = 0 it can be shown that k ◊ H = 0.
     
Hence, the wave vector k is perpendicular  to
 H. Therefore,
   relations k ◊ E = 0 and k ◊ H = 0 indicate
the
that the electromagnetic
 field vectors E and H (or B, as B = m0H) both are perpendicular to the direction of
propagation vector k . It means that the electromagnetic waves are transverse in nature.
11.25 maxweLL’s equatiOns in isOtrOpiC

LO11
DieLeCtriC meDium: em waVe prOpOgatiOn

In an isotropic dielectric medium, the conduction or free current density J and volume charge density r are
      
zero. Further, the displacement vector D and the magnetic field B are defined as D = eE and B = mH. In fact D
          
= e0E + P = e0E + e0ceE ∫ eE together with e = e0(1 + ce) and B = m0H + m0M = m0H + m0cmH ∫ mH together
with m = m0(1 + cm) for the isotropic linear dielectric (polarizable and magnetic) medium. Here, the vectors
 
P and M give respectively the measure of polarization and magnetization of the medium. However, for the
dielectric medium, it would be sufficient to remember that e0 and m0 of free space have been simply replaced
with e and m. Hence, for dielectric medium
 
J = 0 (or s = 0, f = 0, D = eE and B = mH)
where e and m, which are respectively the absolute permittivity and permeability of the medium. Under this
situation, we can express the Maxwell’s equation as
 
—◊ E = 0 (i)
 
—◊ H = 0 (ii)

  ∂H
—¥ E = -m (iii)
∂t

  ∂E
—¥H =e (iv)
∂t
Taking curl of Eq. (iii), we get

    È ∂H ˘
— ¥ (— ¥ E ) = — ¥ Í - m ˙
Î ∂t ˚
or
    ∂  
—(—◊ E ) - —2 E = - m (— ¥ H )
∂t

2
∂ E
or 0 - —2 E = - me [Using Eqs (i) and (iv)]
∂t

 2
∂ E
—2 E = me 2 (v)
∂t
Similarly, taking curl of Eq. (iv) and using Eqs. (ii) and (iii), we get

 ∂2 H
2
— H = me 2 (vi)
∂t
1
As discussed earlier, gives the phase velocity of the wave in the medium. If we represent this as v, we
me
obtain from Eqs. (v) and (vi)

 1 ∂2 E
—2 E - 2 2 = 0
v ∂t

and
 1 ∂2 H
2
— H- 2 =0
v ∂t 2
Eqs. (v) and (vi) are the wave equations in an isotropic linear dielectric medium.
1 1
Now, v = = ( m = m0mr and e = e0er)
me m0 m r e 0e r
c È 1 ˘
or v= Í c= (vii)
mr e r Î m0 e 0 ˙˚
Eq. (vii) shows that the propagation velocity of an electromagnetic wave in a dielectric medium is less than
that in free space.
c
Also, refractive index == mr e r (viii)
v
For non-magnetic dielectric medium mr ª 1. Hence, refractive index = er or Refractive index
= Relative permittivity
11.26 maxweLL’s equatiOns in COnDuCting

LO11
meDium: em waVe prOpagatiOn anD skin Depth
We consider a linear and isotropic conducting medium whose permeability is m, permittivity is e and the
conductivity is s. Under this situation, we can write the Maxwell’s equation as
 
—◊ E = 0 (i)
 
—◊ H = 0 (ii)
 
  ∂B ∂H
—¥E=- = -m (iii)
∂t ∂t

   ∂E
—¥H = J +e
∂t

   ∂E  
or — ¥ H = sE + e [ J = sE ] (iv)
∂t
Taking curl of Eq. (iii), we have

    È ∂H ˘
— ¥ (— ¥ E ) = — ¥ Í - m ˙
Î ∂t ˚

È  ∂H ˘
= - m Í— ¥ ˙
Î ∂t ˚
∂  
= - m [— ¥ H ]
∂t

∂ È  ∂E ˘
= - m Ís E + e ˙ [using Eq. (iv)]
∂t Î ∂t ˚
 
   ∂E ∂2 E
or — ¥ (— ¥ E ) = - ms - me 2
∂t ∂t
           
Also, — ¥ (— ¥ E ) = — (— ◊ E ) – (— ◊ —)E = –—2E [\— ◊ E = 0 from Eq. (i)]
 
 ∂E ∂2 E
\ 2
- — E = - ms - me 2
∂t ∂t
 
 ∂E ∂ E2
or 2
— E = + ms + me 2 (v)
∂t ∂t

Eq. (v) is the electromagnetic wave equation for the electric field E in a conducting medium.
In case of non-conducting medium s = 0. Hence, from Eq. (v)

 ∂2 E
—2 E = me 2 (vi)
∂t

∂E
Eqs (v) and (vi) show that the term ms is the dissipative term which allows the current to flow through
∂t
the medium due to the appearance of conductivity s.
Now, by taking curl of Eq. (iv), we obtain


    È  ∂E ˘
— ¥ (— ¥ H ) = — ¥ Ís E + e ˙
Î ∂t ˚
  ∂  
= s (— ¥ E ) + e (— ¥ E )
∂t
 
Ê ∂H ˆ ∂Ê ∂H ˆ [Using Eq. (iii)]
= s Á -m ˜ + e ÁË - m ˜
Ë ∂t ¯ ∂t ∂t ¯
 
∂H ∂2 H
= - ms - em 2
∂t ∂t
        
 — ¥ (— ¥ H) = —(— ◊ H) – (— ◊ —)H
  
= 0 – —2H [\ — ◊ H = 0 from Eq. (ii)]
 
 ∂H ∂2 H
\
2
-— H = - ms - me 2
∂t ∂t
 
or
 ∂H 2
∂ H (vii)
— 2 H = ms + me 2
∂t ∂t 
Equations (vii) represents the electromagnetic wave equation for magnetic field (H) in conducting medium.
In case of non-conducting medium s = 0, the wave equation takes the form

 ∂2 H (viii)
2
— H = me 2
∂t 
∂H
Equations (vii) show that the term ms is the dissipative term which allows the current to flow through
the conducting medium. ∂t
11.26.1 solution of Wave equation

Equations (v) and (vii) are called inhomogeneous wave equation due to the presence of dissipative term.
These equations in one-dimension (along z-axis) are written as
  
∂2 E ∂E ∂2 E
= ms + me 2 (ix)
∂z 2 ∂t ∂t
  
∂2 H ∂H ∂2 H
= ms + me (x)
∂z 2 ∂t ∂t 2
We assume the following plane wave solutions to the above equations
 
E (z, t) = E0 ei(kz – wt) (xi)
  i(kz – wt)
and H (z, t) = H 0 e (xii)
The use of Eq. (xi) in Eq. (ix) and Eq. (xii) in Eq. (x) leads
k2 = mew2 + imsw (xiii)
This relation is called dispersion relation that governs the electromagnetic wave propagation in a conducting
medium. Equation (xiii) suggests that the quantity k, i.e., wave number, will be a complex quantity. So, we
assume
k = kr + iki (xiv)
 
With this the fields E and H become
 
E ( z , t ) = E0e- ki z ◊ ei ( kr z - w t ) (xv)
 
and H ( z , t ) = H 0e- ki z ei ( kr z - wt ) (xvi)
11.26.2 skin Depth

The expressions (xv) and (xvi) follow that the amplitude of the electric field E is E0 e ki z and that of the

magnetic field H is H 0 e ki z . Hence the amplitude of electromagnetic wave will decrease exponentially as it
propagates through the conductor. This is called the attenuation of the wave and the distance through which
the amplitude is reduced by a factor of 1/e is called the skin depth or penetration depth d. At z = d, the am-
plitude is E0/e. Hence
E0 e kid = E0 /e (i)
This gives the skin depth as
1
d= (ii)
ki
Equation (ii) shows that the imaginary part of the wave number k is the measure of the skin depth. However,
the real part kr of k determines the wave propagation characteristics in the following manner.
Wavelength l = 2p/kr (iii)
Phase velocity n = w/kr (iv)
c ckr
Refractive index n = = (v)
n w
By putting k = kr + i ki in Eq. (xiii) we obtain
me ÈÍ ˘
2
Ês ˆ
kr = w 1 + Á ˜ + 1˙ (vi)
2 ÍÎ Ë we ¯ ˙˚
me ÈÍ ˘
2
Ês ˆ
and ki = w 1 + Á ˜ - 1˙ (vii)
2 ÍÎ Ë we ¯ ˙˚
For good conductors, s >> we. This condition when put in Eqs. (vi) and (vii) gives
me s
kr = ki = w
2 we
wms
fi kr = ki = = p f ms (viii)
2
Hence, the skin depth is given by
1
d= (ix)
p f ms
Since d is inversely proportional to f, which is the frequency of electromagnetic wave, high frequency waves
are found to penetrate less into the conductor. Also, the penetration will be less in the medium having high
conductivity s. Ideally an electromagnetic wave will not penetrate into a perfect conductor as s = •.
 
11.26.3 Phase relationship of e and B Fields
In view of the imaginary wave number ki, we can also make another observation with regard to the phase
  
difference between E and H vectors. If we take the direction of E field along the x-axis, then

  ∂H
—¥ E =-m gives
∂t
 kE
H ( z , t ) = ˆj 0 e - ki z ei ( kr z - wt ) (x)
wm
Ê k¢ˆ
Since k is the complex quantity, it can be represented as k = k ¢eiq k . Here k ¢ = kr2 + ki2 and q k = tan -1 Á i ˜ .
 Ë kr ¯
Then the expression of H becomes
 k ¢E0 - ki z i ( kr z - wt + q k )
H ( z , t ) = ˆj e e (xi)
wm

A comparison of Eq. (xi) with E ( z , t ) = iE ˆ 0e - ki z ei ( kr z - wt ) reveals that the electric field and magnetic field
vectors do not remain in phase when electromagnetic wave propagates in a conducing medium. This is in
contrast to the cases of vacuum and dielectrics.
11.27 eLeCtrOmagnetiC energy Density LO11

It can be proved that the work done in assembling a static charge distribution (number n) against the Coulomb
repulsion of like charges is
n
1 
WE =
2
Â q jV ( r j ) (i)
j =1
For a volume charge density r, this equation takes the form

1 (ii)
2Ú
WE = rVdX
Here
 dX is the volume element.
 Now
 the above equation can be written in terms of
 the resulting electric field
E if we apply Gauss’s law e0— ◊ E = r and mention the potential V in terms of E . This yields the following
relation where the integration is over all the space containing the whole charge distribution.
e0
2 Ú
WE = E 2 dX (iii)
The same way we can derive an expression for the work done on a unit charge against the back emf in one
trip around the circuit, as follows
1
2m0 Ú
WB = B 2 dX (iv)
Here B is the resulting magnetic field. Eqs. (iii) and (iv) suggest that the total energy stored in electromagnetic
field would be
1 Ê B2 ˆ
2 Ú ÁË
WEM = e E 2
+ dX (v)
m0 ˜¯
0
Therefore, the electromagnetic energy density can be obtained as
1Ê B2 ˆ (vi)
U EM = Á e0 E 2 +
2Ë m0 ˜¯
For a monochromatic plane electromagnetic wave, B =

E Ê 1 ˆ is the speed of light. Hence,
, where c ∫
c ÁË m0 e 0 ˜¯

it can be seen that the contribution of magnetic field B to the energy density is the same as that of the electric
field E .
11.28 pOynting VeCtOr anD pOynting theOrem LO12

The electromagnetic waves carry energy when they propagate and there is an energy density associated with
both the electric and magnetic fields.
The amount of energy flowing through unit area, perpendicular to the direction of energy propagation per
unit time, i.e., the rate of energy transport per unit area,  is called
 the poynting vector. It is also termed as
instantaneous energy flux density and is represented by S (or P , sometimes). Mathematically it is defined
  
  S =E¥H
where E and H represent the instantaneous values of the electric  and magnetic field vectors. This is clear that
the rate of  energy
 transport S is perpendicular
 to both E and H and is in the direction of propagation of the
wave, as E ¥ H is in the direction of k . Since the poynting vector represents the rate of energy transport per
unit area, its units are W / m2.
Derivation: We can calculate the energy density carried by electromagnetic waves with the help of Maxwell’s
equations given below.

div D = 0 (i)

div B = 0 (ii)

 - ∂B
curl E = (iii)
∂t

  ∂D
curl H = J + (iv)
∂t
 
Take scalar (dot) product of Eq. (iii) and Eq. (iv) with H and E respectively, i.e.,

   ∂B
H ◊ curl E = - H ◊ (v)
∂t

     ∂D
and E ◊ curl H = E ◊ J + E ◊ (vi)
∂t
Subtract Eq. (vi) from Eq. (v), i.e.,
 
     ∂B    ∂D
H ◊ curl E - E ◊ curl H = - H ◊ - E◊J - E◊
∂t ∂t
 
Ê  ∂B  ∂D ˆ  
= -ÁH ◊ + E◊ ˜ - E◊J
Ë ∂t ∂t ¯
 
or   Ê  ∂B  ∂D ˆ   (vii)
div (E ¥ H ) = - Á H ◊ + E◊ ˜ - E◊J
Ë ∂t ∂t ¯
     
[ div ( A ¥ B ) = B ◊ curl A - A ◊ curl B ]
   
Using the relations B = m H and D = e E , we can get

 ∂D  ∂  1 ∂ ∂ Ê 1  ˆ
E◊ = E ◊ (e E ) = e ( E ) 2 = Á E ◊ D˜  
∂t ∂t 2 ∂t ∂t Ë 2 ¯ [ E 2 = E ◊ E ]

 ∂B  ∂  1 ∂ ∂ Ê 1  ˆ  
H◊ = H ◊ ( m H ) = m ( H ) 2 = Á H ◊ B˜ [ H 2 = H ◊ H ]
∂t ∂t 2 ∂t ∂t Ë 2 ¯
Now Eq. (vii) can be written as
  ∂ È1     ˘  
div ( E ¥ H ) = Í ( H ◊ B + E ◊ D) ˙ - E ◊ J
∂t Î 2 ˚
  ∂ È1     ˘  
or E ◊ J = Í ( H ◊ B + E ◊ D) ˙ - div ( E ¥ H ) (viii)
∂t Î 2 ˚
Integrating Eq. (viii) over a volume V enclosed by a surface S, we get
v
 
Ú E ◊ JdV = - Ú ÍÎ ∂t
v
2 {
È∂ 1     ˘
}  
( H ◊ B + E ◊ D) ˙ dV - Ú div ( E ¥ H ) dV
˚ v
  ∂ ÈÊ 1 1 ˆ˘   
Ú E ◊ JdV = - ∂t Ú ÍÎÁË 2 m H Ú
2
or + e E 2 ˜ ˙ dV -  ( E ¥ H ) ◊ dS (ix)
2 ¯ ˚
v v s
        
[ B = m H , D = e E and Ú div ( E ¥ H )dV = 
Ú ( E ¥ H ) ◊ dS ]
v s
Eq. (ix) can also be written as
  ∂ È1 1 2˘   
Ú E ◊ J dV = - ∂t Ú ÍÎ 2 m H + 2 e E ˙˚ dV - Ú ( E ¥ H ) ◊ dS
2 (x)
v v s
Interpretation
 
(a) Ú ( E ◊ J ) dV : This term represents the rate of energy transferred into the electromagnetic field
v
through the motion of charges in the volume V, i.e., the total power dissipated in a volume V.
∂ È1 1 ˘ 1 1
(b) Ú Í
∂t v Î 2
m H 2 + e E 2 ˙ dV : The terms m H 2 and e E 2 represent the energy stored in electric and
2 ˚ 2 2
magnetic fields respectively and their sum will be equal to the total energy stored in electromagnetic
field. Therefore, this total expression represents the rate of decrease of energy stored in volume V
due to electric and magnetic fields.
  
(c) Ú ( E ¥ H ) ◊ dS : This term represents the amount of electromagnetic energy crossing the closed
s
surface per second

  or the rate of flow of outward energy
 through
  the
 surface S enclosing volume V.
The vector (E ¥ H) is known as the poynting vector S or S = (E ¥ H).
Equation (x) is also known as poynting theorem or work-energy theorem, according to which the power
transferred into the electromagnetic field is equal to the sum of the time rate of change of electromagnetic
energy within a certain volume and the time rate of the energy flowing out through the boundary surface. This
is also called as the energy conservation law in electromagnetism.
11.29 waVe prOpagatiOn in BOunDeD system: waVeguiDe LO12

The first waveguide was proposed by Thomson in 1893 and was experimentally verified by Lodge in 1894. A
structure that can guide waves like electromagnetic waves, light waves or sound waves is called waveguide.
For each type of the wave there are different types of waveguides. For example, depending on the frequency
of electromagnetic wave the waveguide can be constructed from either conductive or dielectric material. Such
electromagnetic waveguides are especially useful in the microwave and optical frequency ranges. The
waveguides used at optical frequencies (optical waveguides) are typically dielectric waveguides. In such
waveguides, a dielectric material with high permittivity (i.e., high index of refraction) is surrounded by
a material with lower permittivity. This type of structure guides optical waves by the process of total
internal reflection. The most common optical waveguide is optical fibre. On the other hand, a structure
that guides sound waves is called an acoustic waveguide. A duct for sound propagation also behaves as a
transmission line. The duct contains some medium like air which supports the propagation of sound waves.
11.29.1 electromagnetic Waveguides
The original and the most common meaning of waveguide is a hollow metal pipe used for guiding the waves.
The electromagnetic waves in such waveguides may be imagined as waves travelling down the guide in a
zig zag path as these waves are repeatedly reflected between opposite walls of the guide (for example,
rectangular waveguide). The first mathematical analysis of the propagating modes (waves) within a hollow
metal cylinder was performed by Rayleigh in 1897.
To function properly, a waveguide must have a certain minimum diameter relative to the wavelength of
the signal. If the waveguide is too narrow or the frequency is too low (the wavelength is too long), the
electromagnetic field cannot propagate. There is a minimum frequency, known as cutoff frequency for the
propagation of the wave, i.e., a wave can propagate only if its frequency is larger than the cutoff frequency.
The cutoff frequency is decided by the dimensions of the waveguide.
11.29.2 Modes in Waveguides
In order to analyse the mode (wave) propagation in the waveguide, we solve the Maxwell’s equations
along with appropriate bounding conditions determined by the properties of the materials and their
interfaces. These equations admit multiple solutions, or modes, which are origin functions of the
equation system. The propagation of the waveguide modes depends on the operating wavelength and
polarization, and shape and size of the waveguide. The longitudinal mode can be realised in a cavity
(closed end waveguide), the longitudinal mode is particularly standing wave pattern formed by the
waves confined in the cavity. However, a number of transverse modes can be excited in the waveguide,
which are classified below.
(1) Transverse Electric (TE) Modes: These modes do not have electric field in the direction of
propagation. So electric field vector is in transverse direction.
(2) Transverse Magnetic (TM) Modes: These modes have no magnetic field in the direction of
propagation. So magnetic field vector is in transverse direction.
(3) Transverse Electromagnetic (TEM) Modes: These modes have no electric and magnetic fields in
the direction of mode propagation. In hollow waveguides, TEM modes are not possible because as
per Maxwell’s equation the electric field then must have zero divergence, zero curl and be zero at
the boundaries. This will result in a zero field or —2E = 0. However, TEM modes can propagate in
a coaxial cable.
(4) Hybrid Modes: These modes have both electric and magnetic field components in the direction of
propagation. The mode for which the cutoff frequency is the minimum is called the fundamental
mode. For example, Transverse Electric TE10 mode is the fundamental mode for rectangular
waveguide whereas TE11 mode is the fundamental mode for circular waveguide.
11.30 COaxiaL CaBLe LO12

A transmission line is the material medium that forms a path from one place to another for directing the
transmission of energy like electromagnetic waves, acoustic waves or electric power. Open wire transmission
lines have the property that the electromagnetic wave propagating down the line extends into the space
surrounding the parallel wires. The transmission lines cannot be bent, twisted or otherwise shaped without
changing their characteristic impedance. They also cannot be attached to anything conductive, as the extended
fields will induce currents in the nearby conductors. This will cause unwanted radiation and detuning of the
line. However, this problem is solved by coaxial lines that confine the electromagnetic wave to the area
inside the cable, as the coaxial cable consists of a round conducting wire surrounded by an insulating spaces
(called the dielectric), surrounded by a cylindrical conducting sheath that is usually surrounded by a final
insulating layer (called jacket). Here the transmission of energy occurs totally through the dielectric inside the
cable between the conductors. Coaxial lines (cables) can therefore be bent and moderately twisted without
negative effects. Also, they can be strapped to conductive supports without inducing unwanted currents in
them. Coaxial cable is used as a high-frequency transmission line to carry a high frequency or broadband
signals. Coaxial cables may be flexible or rigid depending on the type of sheath. Rigid types have a solid
sheath, whereas flexible types have an interwoven sheath usually of thin copper wire. The inner insulator
has a significant effect on the cable’s properties such as its characteristic impedance and its attenuation. The
characteristic impedance in ohms is calculated from the ratio of the inner diameter d and outer diameter D of
the dielectric and the dielectric constant er as below
138 Ê Dˆ
Z0 = log Á ˜
er Ëd¯
The most common impedances that are widely used are 50 or 52 W for industrial and commercial radio
frequency applications, and 75 W for domestic television and radio, although other impedances are available
for specified applications.
11.30.1 signal Propagation in coaxial cables

In radio-frequency applications up to a few GHz (109Hz) the wave propagates only in the transverse
electromagnetic (TEM) mode which means the electric and magnetic fields are both perpendicular to the
direction of propagation. However, above a certain cutoff frequency, transverse electric (TE) and / or transverse
magnetic (TM) modes can also propagate, as they do in a waveguide. It is usually undesirable to transmit signals
above the cutoff frequency, since it may cause multiple modes with different phase velocities to propagate,
interfering with each other. This cutoff frequency is roughly inversely proportional to the outer diameter.
11.30.2 uses of coaxial cables

Short length coaxial cables are commonly used to connect home video equipment and in measurement
electronics. They are used to be common for implementing computer networks (particularly ethernet) but
twisted pair cables have replaced them in most applications except in the growing consumer cable modem
market for broadband internet access.
Long distance coaxial cable is used to connect radio networks and television networks, though this has largely
been superseded by other more high-tech methods which use fibre optics, satellite etc. Micro coaxial cables
are used in a range of consumer devices, military, equipment and also in ultrasound scanning equipment.
s UmmarY

✦ The separation of charges produces electric field, whereas the motion of charges generates current and
hence the magnetic field. When these fields are time varying they are coupled with each other through
the Maxwell’s equations. With the help of the Maxwell’s equations, we can derive wave equation,
based on which the propagation of electromagnetic waves can be investigated in different media.
✦ If a charge is distributed continuously in a medium, it can be expressed in terms of a physical quantity
known as charge density. There are three types of charge densities, namely linear charge density l,
surface charge density s and volume charge density r.
✦ Gauss’s law of electricity states that the net outward electric flux through any closed surface drawn in
an electric field is equal to 1 / e0 times the total charge enclosed within the surface. Here it is clear that
the charge located outside the surface does not contribute to the electric flux and hence the electric field.

✦ The del operator is the differential operator, which is represented by — and is given by
 ∂ ∂ ∂
— = iˆ + ˆj + kˆ . It is not a vector in itself, but when it operates on a scalar function it provides
∂x ∂y ∂z
 
the resultant as a vector. —F does not mean a multiplication of — with F rather it is an instruction to
differentiate. Here we should say that — is a vector operator that acts upon F.
✦ The gradient of a scalar field F is a vector quantity which represents both the magnitude and the
direction of the maximum space rate of increase of F. The gradient of a vector field cannot be realised.

✦ The divergence of a vector field A is defined as the net outward flux per unit volume over a closed
surface. The gradient at a given point is a measure of how much the vector A spreads out, i.e., diverges,
from that point.

✦ The
 curl of a vector field A is a rotational vector. Its magnitude would be the maximum circulation of
A per unit area. Its direction is the normal
 direction of the area when the area is oriented so as to make
the circulation maximum. The curl of A at some point O is a measure of how much the vector A curls
around the point O. a2
For a function g(x), the fundamental theorem of calculus states that Ú ÊÁ ˆ˜ dx = g (a2 ) - g (a1 ) It
dg
✦
Ë dx ¯
a1
means the total change in the function g(x) can be obtained by simply subtracting the values of the
function at the points a2 and a1.
a2
 
✦ For a scalar function F, Ú —F ◊ dl = F (a2 ) - F (a1 ) . This is called the fundamental theorem for gradient,
a1
according to which the integral of a derivative over end points a1 and a2 is given by the value of the
function at the boundaries, i.e., the points a1 and a2.
   
✦ According to Gauss’s or Green’s theorem, Ú F )dV = Ú F ◊ dS . Here dV is the volume element and
( —◊
V S

dS is the surface element. This theorem states that the integral of a derivative (here the divergence) over
a region (here the volume) is equal to the value of the function at the boundary (here the surface). Since
the boundary of a volume is always a closed surface, the R.H.S. is the integral over closed surface.
Evidently this theorem converts the volume integral into the surface integral.
✦ Stokes’ theorem states that the integral of the curl of a vector function over a patch of surface is equal
to the value of the function at the perimeter of the patch. So here the derivative is the curl, region is the
surface and the boundary is the perimeter of the patch of the surface. Therefore,
    
Ú (— ¥ F ) ◊ dS = Ú F ◊ dl
S C
Clearly, this theorem converts the surface integral into the line integral.

✦ If we represent the charge density by r and the current density by J , then the continuity equation is
  ∂r
written as —◊ J + = 0 . In this equation, the second term is the time rate of change of charge density,
∂t 
i.e., the current. The first term represents the divergence of J , i.e., the measure of spreading of the
current through a surface. This spreading is balanced by the time rate of change of charge density.
Therefore, this equation just tells us the conservation of charges.
✦ When the charges are in motion, the electric and magnetic fields are associated with this motion, which
will have variations in both the space and time. These electric and magnetic fields are interrelated.
This phenomenon is called electromagnetism, which is summarised by the set of equations, known as
Maxwell’s equations. The Maxwell’s equations are nothing but  are the representation
 of the basic laws
of electromagnetism. If the electric field is represented by E , magnetic field by B, current density by
 
       ∂B    ∂D
J , then the Maxwell’s equations are , —◊ D = r , —◊ B = 0, — ¥ E = - and — ¥ H = J + . Here
    ∂t ∂t
B = m H and D = e E for the linear media.
 
✦ The Maxwell’s first equation — ◊ D = r represents the Gauss’s Law for electricity along with r as the
free charge density.
✦ The area integral of a vector field determines the net source of the field (function). Since the Maxwell’s
   
second equation — ◊ B = 0 can be written as  Ú B ◊ dS = 0 , it says that the net magnetic flux out of any
closed surface is zero. This is because the magnetic flux directed inward toward the south pole, of a
magnetic dipole kept in any closed surface, will be equal to the flux outward the north pole. Therefore,
the net flux is zero for dipole sources. For a magnetic monopole source, the value of the area integral
 
Ú B ◊ dS would be finite. Since the divergence of a vector field is proportional to the density of point
source, the Maxwell’s second equation simply says that there are no magnetic monopoles.

  ∂B
✦ The Maxwell’s third equation — ¥ E = - when written in the integral form states that the line
∂t
integral of the electric field around a closed loop is equal to the negative rate of change of the magnetic
flux through the area enclosed by the loop. The line integral basically is the generated voltage or emf
in the loop. Therefore, the physical interpretation of Maxwell’s third equation is that the changing
magnetic field induces electric field.
 
✦ For static electric field E , i.e., when E does not change with time , the second term of the R.H.S. of the

   ∂D
Maxwell’s fourth equation — ¥ H = J + vanishes and then the integral form of this equation says
∂t
that the line integral of the magnetic field around a closed loop is proportional to the electric current
flowing through the loop. This form of the Maxwell’s equation is useful for calculating the magnetic
field for simple geometries. However, this equation more specifically reveals that the changing

∂D
electric field induces magnetic field as is nothing but the rate of change of electric field. This
∂t
is complimentary to the meaning of the Maxwell’s third equation. Therefore, they together form the
electromagnetic fields or electromagnetic waves, where both electric and magnetic fields propagate
together and the change in one field induces the other field.

  ∂D
✦ The vector D(= eE ) is called the displacement vector and the term is called the displacement
∂t
current. The displacement current is postulated in a dielectric when electric stress or potential gradient
is varied. It is different from a normal or conduction current, as it is not accompanied by the motion of
current carriers in the dielectrics. The concept of introducing this term in the Ampere’s law was given
by Maxwell for the completion of his electromagnetic equations, as the Ampere’s law is valid only for
the steady state or static fields. In order to make it consistent with the time varying fields, i.e., also for
satisfying the continuity equation, this additional term is required.
✦ Electromagnetic fields, i.e., the electric field and magnetic field are transverse in nature that always
remain perpendicular to each other and also in the transverse direction to the direction of wave
propagation.
✦ An electromagnetic wave of frequency w and wave number k satisfies the dispersion relation
w
k2 – w2m0e0 = 0 in free space or vacuum. The wave propagates at the phase velocity vp = = c, the
k
speed of light, in this medium.
✦ An electromagnetic wave of frequency w and wave number k satisfies the dispersion relation
k2 – w2me = 0 in a dielectric medium whose permittivity is e and permeability is m. Its phase velocity is
w 1 c
vp = = = . It means this wave propagates at slower speed in dielectric medium compared
k me mr e r
with the vacuum or free space.
✦ For non-magnetic dielectric medium, mr = 1 and its refractive index = relative permittivity .
✦ The wave equation for an electromagnetic wave in a conducting medium is modified due to the term
 
∂E ∂H
ms or ms . This term is called dissipative term which allows the current to flow through the
∂t ∂t
medium due to the appearance of conductivity s.
✦ The wave number k in the case of the propagation of an electromagnetic wave in a conducting medium
does not remain a real quantity rather becomes a complex quantity due to the development of its
imaginary part. The imaginary wave number leads to the attenuation of the wave in the medium, i.e.,
exponential decay of the amplitude/field of the wave when it enters the conductors.
✦ The distance through which the amplitude of the wave is reduced by a factor of 1/e is called the skin
depth or penetration depth d, given by d = 1/ p f ms , where f is the frequency of the wave, m is the
permeability and s is the conductivity of the conducting medium.
✦ Unlike the cases of vacuum/free space and dielectric medium, the electric field and magnetic field
vectors do not remain in phase when electromagnetic wave propagates in a conducting medium. This
happens due to the presence of imaginary part of the wave number.
✦ The real part kr of the wave number k determines the wave propagation characteristics when an
electromagnetic wave propagates in a conducting medium.
For example,
Phase velocity vp = w /kr
Wavelength l = 2p/kr
c ckr
Refractive index n = =
vp w
✦ The electromagnetic waves carry energy when they propagate and there is an energy density associated
with both the electric and magnetic fields. The amount of energy flowing through unit area, perpendicular
to the direction of energy propagation per unit time, i.e., the rate of energy transport per unit area, is
called the poynting vector. It is also known as instantaneous energy flux density and is represented by
  
S = E ¥ H.
✦ There are structures that can guide waves like electromagnetic waves, light waves or sound waves. These
structures are called waveguides. For each type of the wave there are different types of waveguides.
For example, depending on the frequency of electromagnetic wave the waveguide can be constructed
from either conductive or dielectric material. Such electromagnetic waveguides are especially useful
in the microwave and optical frequency ranges. The waveguides used at optical frequencies (optical
waveguides) are typically dielectric waveguides.
✦ The waveguides support different types of wave propagation, which are called modes. In order to
analyze the mode propagation, we solve the Maxwell’s equations along with appropriate boundary
conditions determined by the properties of the materials and their interfaces. These equations admit
multiple solutions, or modes, which are origin functions of the equation system. The mode propagation
in the waveguides depends on the operating wavelength and polarization, and shape and size of the
waveguide.
✦ A number of transverse modes can be excited in the waveguide. For example, transverse electric (TE)
modes which do not have electric field in the direction of propagation, transverse magnetic (TM) modes
which have no magnetic field in the direction of propagation, transverse electromagnetic (TEM) modes
which have no electric and magnetic fields in the direction of mode propagation. In hollow waveguides,
TEM modes are not possible. However, TEM modes can propagate in a coaxial cable. Also hybrid modes
can be excited which have both electric and magnetic field components in the direction of propagation.
✦ The mode for which the cutoff frequency is the minimum is called the fundamental mode. For example,
TE10 mode is the fundamental mode for rectangular waveguide whereas TE11 mode is the fundamental
mode for circular waveguide.
✦ A transmission line forms a path from one place to another for directing the transmission of energy
like electromagnetic waves, acoustic waves or electric power. The transmission lines cannot be
bent, twisted or otherwise shaped without changing their characteristic impedance. They cannot be
attached to anything conductive, as the extended fields will induce currents in the nearby conductors.
This will cause unwanted radiation and detuning of the line. However, coaxial lines that confine the
electromagnetic wave to the area inside the cable solve this problem. Here the transmission of energy
occurs totally through the dielectric inside the cable between the conductors. Coaxial lines (cables)
can therefore be bent and moderately twisted without negative effects. Also, they can be strapped to
conductive supports without inducing unwanted currents in them.
✦ In radio frequency applications up to a few GHz (109Hz) the wave propagates only in the form of
transverse electromagnetic (TEM) mode. However, above a certain cutoff frequency, transverse electric
(TE) and / or transverse magnetic (TM) modes can also propagate, as they do in a waveguide. It is
usually undesirable to transmit signals above the cutoff frequency, since it may cause multiple modes
with different phase velocities to propagate, interfering with each other.
s olVeD e XamPles
 
E xamplE 1 Find the value of —r n where r = xiˆ + yjˆ + zkˆ .
Solution
 Ê ∂ ∂ ∂ˆ
—r n = Á iˆ + ˆj + kˆ ˜ ( x 2 + y 2 + z 2 ) n /2
Ë ∂x ∂y ∂z ¯
where r = (x2 + y2 + z2)1/2
 Ï ¸ Ï ¸
{ }
n n
-1 -1 n n
—r n = Ì n ( x 2 + y 2 + z 2 ) 2 2 x ˝ iˆ + Ì n ( x 2 + y 2 + z 2 ) 2 2 y ˝ ˆj + ( x 2 + y 2 + z 2 ) - 1 kˆ
Ó2 ˛ Ó2 ˛ 2 2
Ên ˆ
Á - 1˜¯
= n ( x 2 + y 2 + z 2 )Ë 2 ˆ + iy
(ix ˆ)
ˆ + kz

= nr n - 2 r
È1˘ n 
E xamplE 2 Prove that — Í n ˙ = - n + 2 r , where r = x2 + y 2 + z 2 .
Îr ˚ r
Ê 1ˆ Ê ∂ ∂ ∂ˆ
Solution — Á n ˜ = Á iˆ + ˆj + kˆ ˜ r - n
Ë r ¯ Ë ∂x ∂y ∂z ¯
-n
Ê ∂ ∂ ∂ˆ
= Á iˆ + ˆj + kˆ ˜ ( x 2 + y 2 + z 2 ) 2
Ë ∂x ∂y ∂z ¯
Ï -n ¸ Ï -n ¸ Ï n ¸
= Ì - n ( x 2 + y 2 + z 2 ) 2 2 x ˝ iˆ + Ì - n ( x 2 + y 2 + z 2 ) 2 2 y ˝ ˆj + Ì - n ( x 2 + y 2 + z 2 ) 2 ˝ kˆ
-1 -1 -1
Ó 2 ˛ Ó 2 ˛ Ó 2 ˛
Ê n+2 ˆ
-Á
Ë 2 ¯˜ ˆ)
= - n( x 2 + y 2 + z 2 ) ˆ + ˆjy + kz
(ix
-n 
= n+2 r
r
Hence proved.

E xamplE 3 If f = x3/2 + y3/2 + z3/2, find — f.
 Ê ∂ ∂ ∂ˆ
Solution —f = Á iˆ + ˆj + kˆ ˜ ( x3/ 2 + y 3/ 2 + z 3/ 2 )
Ë ∂x ∂y ∂z ¯
3 ˆ 1/ 2 ˆ 1/ 2 ˆ 1/ 2
= (ix + jy + kz )
2
E xamplE 4 If f(x, y, z) = 3x2y – yz2, find grad f at point (1, 2, –1).
Ê ∂ ∂ ∂ˆ
Solution grad f = —f = Á iˆ + ˆj + kˆ ˜ (3 x 2 y - yz 2 )
Ë ∂x ∂y ∂z ¯
= iˆ(6 xy ) + ˆj (3 x 2 - z 2 ) + kˆ ( -2 yz )
Value of grad f at point (1, 2, –1) = 12iˆ + 2 ˆj + 4kˆ

 
E xamplE 5 If f is a scalar field and A is vector field find the value of div (fA).
  
Solution div (fA) = —◊ (fA)
Ê ∂ ∂ ∂ˆ
= Á iˆ + ˆj + kˆ ˜ ◊ (f Axiˆ + f Ay ˆj + f Az kˆ)
Ë ∂x ∂y ∂x ¯
∂ ∂ ∂
= (f Ax ) + (f Ay ) + (f Az )
∂x ∂y ∂z
∂f ∂A ∂f ∂Ay ∂f ∂A
= Ax + f x + Ay + f + Az + f z
∂x ∂x ∂y ∂y ∂z ∂z
Ê ∂f ∂f ∂f ˆ Ê ∂A ∂Ay ∂Az ˆ
=Á Ax + Ay + Az + f Á x + +
Ë ∂x ∂y ∂z ˜¯ Ë ∂x ∂y ∂z ˜¯
∂A
Ê ∂f
= Á iˆ +
∂f ˆ ∂f ˆ ˆ ˆ
j+ ˆ ) + f Ê ∂Ax + y + ∂Az ˆ
k ˜ ◊ (iAx + ˆjAy + kA
Ë ∂x ∂y ∂z ¯
z Á
Ë ∂x ∂y ∂z ˜¯
   
div (f A) = —f ◊ A + f (—◊ A)
      
E xamplE 6 Prove that — ◊ (A + B) = — ◊ A + — ◊ B where A and B are differentiable vector functions.
 
Solution A = iA ˆ , B = iB
ˆ x + ˆjAy + kA ˆ
ˆ x + ˆjB y + kB
z z
 
and ( A + B ) = iˆ( Ax + Bx ) + ˆj ( Ay + B y ) + kˆ ( Az + Bz ) (i)
Taking divergence on both sides of Eq. (i), we have
   Ê ∂ ∂ ∂ˆ
—◊ ( A + B ) = Á iˆ + ˆj + kˆ ˜ ◊ [( Ax + Bx )iˆ + ( Ay + B y ) ˆj + ( Az + Bz )kˆ ]
Ë ∂x ∂y ∂z ¯
∂ ∂ ∂
= ( Ax + Bx ) + ( Ay + B y ) + ( Az + Bz )
∂x ∂y ∂z
Ê ∂A ∂Ay ∂Az ˆ Ê ∂Bx ∂B y ∂Bz ˆ
=Á x + + + + +
Ë ∂x ∂y ∂z ˜¯ ÁË ∂x ∂y ∂z ˜¯
      
—◊ ( A + B ) = —◊ A + —◊ B

E xamplE 7 Prove that A = 3 y 2 z 2 iˆ + 3 x 2 z 2 ˆj + 3 x 2 y 2 kˆ is a solenoidal vector.

Solution A = 3 y 2 z 2iˆ + 3x 2 z 2 ˆj + 3x 2 y 2 kˆ
  Ê ∂ ∂ ∂ˆ
—◊ A = Á iˆ + ˆj + kˆ ˜ ◊ (3 y 2 z 2iˆ + 3x 2 z 2 ˆj + 3x 2 y 2 kˆ)
Ë ∂x ∂y ∂z ¯
∂ ∂ ∂
= (3 y 2 z 2 ) + (3x 2 z 2 ) + (3x 2 y 2 )
∂x ∂y ∂z
=0
 
As the divergence of a vector field A is zero, the vector A is solenoidal.

E xamplE 8 Find the constant a, the the vector A = ( x + 3 y )iˆ + (2 y + 3z ) ˆj + ( x + az )kˆ is a solenoidal vector.

Solution A = ( x + 3 y )iˆ + (2 y + 3 z ) ˆj + ( x + az )kˆ
  Ê ∂ ∂ ∂ˆ
div A = —◊ A = Á iˆ + ˆj + kˆ ˜ ◊ [( x + 3 y )iˆ + (2 y + 3z ) ˆj + ( x + az )kˆ ] = 0
Ë ∂x ∂y ∂z ¯
∂ ∂ ∂
or ( x + 3 y ) + (2 y + 3 z ) + ( x + az ) = 0
∂x ∂y ∂z
1+2+a=0
or a = −3
 
 ˆ
E xamplE 9 Calculate the value of —◊ (r 3 r ) where r = (ix ˆ ).
+ ˆjy + kz
 ˆ ˆ
Solution Given r = (ix + ˆjy + kz )
  Ê ∂ ∂ kˆ∂ ˆ ˆ )]
—◊ (r 3r ) = Á iˆ + ˆj + ◊ [( x 2 + y 2 + z 2 )3/ 2 (ix
ˆ + ˆjy + kz
Ë ∂x ∂y ∂z ˜¯
∂ ∂ ∂
= [ x( x 2 + y 2 + z 2 )3/ 2 ] + [ y ( x 2 + y 2 + z 2 )3/ 2 ] + [ z ( x 2 + y 2 + z 2 )3/ 2 ]
∂x ∂y ∂z
∂ 3
Now [ x( x 2 + y 2 + z 2 )3/ 2 ] = ( x 2 + y 2 + z 2 )3/ 2 + x 2 x[ x 2 + y 2 + z 2 ]1/ 2
∂x 2
= r3 + 3x2r
∂
Similarly, [ y ( x 2 + y 2 + z 2 )3/ 2 ] = r 3 + 3 y 2 r
∂y
∂
and [ z ( x 2 + y 2 + z 2 )3/ 2 ] = r 3 + 3z 2 r
∂z
 
so —(r 3r ) = 3r 3 + 3( x 2 + y 2 + z 2 )r
= 3r3 + 3r2 r = 6r3
 -2 z 2 y 2
E xamplE 10 Show that the vector field A = ˆ + z ˆj + 2 yz kˆ is irrotational.
i
x3 x2 x2
 -2 z 2 y z2 2 yz
Solution Given A = 3 iˆ + 2 ˆj + 2 kˆ
x x x

If a vector field A is irrotational, then
  
Curl A = — ¥ A = 0
iˆ ˆj kˆ
  ∂ ∂ ∂
Now — ¥ A =
∂x ∂y ∂z
-2 z 2 y z2 2 yz
x3 x2 x2
È ∂ Ê 2 yz ˆ ∂ Ê z 2 ˆ ˘ È ∂ Ê -2 z 2 y ˆ ∂ Ê 2 yz ˆ ˘ È ∂ Ê z 2 ˆ ∂ Ê -2 z 2 y ˆ ˘
= iˆ Í Á 2 ˜ - Á 2 ˜ ˙ + ˆj Í Á ˜ - Á 2 ˜ ˙ + kˆ Í Á 2 ˜ - Á 3 ˜˙
Î ∂y Ë x ¯ ∂z Ë x ¯ ˚ Î ∂z Ë x ¯ ∂x Ë x ¯ ˚ Î ∂x Ë x ¯ ∂y Ë x ¯ ˚
3
È 2˘
ˆj È -4 yz + 4 yz ˘ + kˆ Í -2 z + 2 z ˙
2
È 2z 2z ˘
= iˆ Í 2 - 2 ˙ + ÍÎ x3
Îx x ˚ x3 ˙˚ Î x3 x3 ˚
 
— ¥ A= 0
  
Hence the curl of a vector field A is zero. So the vector field A is zero, So the vector field A is irrotational.

E xamplE 11 Consider a vector field A = x 2 iˆ + y 2 ˆj + z 2 kˆ
(i) Is the field solenoidal?
(ii) Is the field irrotational?

Solution Given A = x 2iˆ + y 2 ˆj + z 2 kˆ (i)

Now A = iA ˆ
ˆ x + ˆjAy + kA (ii)
z
By using Eqs (i) and (ii), We have
Ax = x2, Ay = y2 and Az = z2
  Ê ∂ ∂ ∂ˆ ˆ
Now — ◊ A = Á iˆ + ˆj + kˆ ˜ . (iA ˆ
x + jAy + kAz )
ˆ
Ë ∂x ∂y ∂z ¯
∂Ax ∂Ay ∂Az
= + +
∂x ∂y ∂z
∂x 2 ∂y 2 ∂z 2
= + + = 2x + 2 y + 2z
∂x ∂y ∂z
 
—◊ Aπ 0

From the above, it is clear that divergence of vector field A is not equal to zero. Hence, this field is not solenoidal.
iˆ ˆj kˆ
  ∂ ∂ ∂ È∂ ∂ ˘ ˆj È ∂ ( x 2 ) - ∂ ( z 2 )˘ + kˆ È ∂ ( y 2 ) - ∂ ( x 2 )˘
(ii) —◊ A= = iˆ Í ( z 2 ) - ( y 2 )˙ + ÍÎ ∂z ˙˚ Í ∂x ˙
∂x ∂y ∂z Î ∂ y ∂ z ˚ ∂x Î ∂y ˚
x2 y2 z2
 
\ —◊ A= 0

Since curl of the vector field A is zero, the field is irrotational.

E xamplE 12 A vector field is given by A = yziˆ + xzjˆ + xykˆ. Show that it is both irrotational and solenoidal.

Solution Given A = yziˆ + xzjˆ + xykˆ

Comparing it with A = Axiˆ + Ay ˆj + Az kˆ
So Ax = yz, Ay = xz and Az = xy

Divergence of vector field A
  Ê ∂ ∂ ∂ˆ ˆ
—◊ A = Á iˆ + ˆj + kˆ ˜ ◊ (iA ˆ
x + jAy + kAz )
ˆ
Ë ∂x ∂y ∂z ¯
∂Ax ∂Ay ∂Az

= + +
∂x ∂y ∂z
∂( yz ) ∂( xz ) ∂( xy )
= + +
∂x ∂y ∂z
 
\ A 0

Hence vector field A is solenoidal.
iˆ ˆj kˆ
   ∂ ∂ ∂
curl A = — ¥ A =
∂x ∂y ∂z
yz xz xy
Ê ∂xy ∂xz ˆ ˆ Ê ∂yz ∂xy ˆ ˆ Ê ∂xz ∂yz ˆ

= iˆ Á - + jÁ - ˜ + kÁ -
Ë ∂y ∂z ¯˜ Ë ∂z ∂x ¯ Ë ∂x ∂y ¯˜
= iˆ( x - x) + ˆj ( y - y ) + kˆ( z - z )
 
\ A 0

Hence vector field A is irrotational.
    
E xamplE 13 Given A = x 2 yiˆ + ( x - y )kˆ . Find (i) —◊ A and (ii) — ¥ A .

Solution Given A = x 2 yiˆ + ( x - y )kˆ

A = Axiˆ + Ay ˆj + Az kˆ
Ax = x2y, Ay = 0 and Az = (x – y)
  Ê ∂ ∂ ∂ˆ ˆ
(i) —◊ A = Á iˆ + ˆj + kˆ ˜ ◊ (iA ˆ
x + jAy + kAz )
ˆ
Ë ∂x ∂y ∂z ¯
∂Ax ∂Ay ∂Az ∂( x 2 y ) ∂(0) ∂( x - y )
= + + = + +
∂x ∂y ∂z ∂x ∂y ∂z
 
—◊ A = 2 xy
iˆ ˆj kˆ
   ∂ ∂ ∂
(ii) curl A = — ¥ A =
∂x ∂y ∂z
x2 y 0 ( x - y)
Ê ∂( x - y ) ∂(0) ˆ ˆ Ê ∂x y ∂( x - y ) ˆ ˆ Ê ∂(0) ∂x y ˆ
2 2
= iˆ Á - ˜ + jÁ - ˜ + kÁ - ˜
Ë ∂y ∂z ¯ Ë ∂z ∂x ¯ Ë ∂x ∂z ¯
= iˆ( -1) + ˆj ( -1) - x 2 kˆ

 
— ¥ A = - iˆ - ˆj - x 2 kˆ

E xamplE 14 Check whether the electrostatic field represented by E = axy 2 ( yiˆ + xjˆ) is conservative or not?
 
Solution If — ¥ E is zero, then the electrostatic field is conservative.

Given E = axy 3iˆ + ax 2 y 2 ˆj

E = iE ˆ
ˆ x + ˆjE y + kE z
So Ex = axy3, Ey = ax2y2 and Ez = 0

iˆ ˆj kˆ
  ∂ ∂ ∂
Then, — ¥ E = π0
∂x ∂y ∂z
axy 3 ax 2 y 2 0
Ê ∂(0) ∂ ˆ Ê ∂ axy3 ∂ ˆ Ê ∂ ( ax 2 y 2 ) ∂ (axy3 ) ˆ
= iˆ Á - (ax 2 y 2)˜ + j Á - (0)˜ + k Á - ˜
Ë ∂y ∂z ¯ ÁË ∂z ∂x ¯˜ ËÁ ∂x ∂y ˜¯
= 2axy 2 - 3axy 2 π 0

Hence E is nonconservative field.
E xamplE 15 If 2000 flux through lines enter a given volume of space and 4000 lines diverge from it, find
the total charge within the volume.
Solution Given f1 = 2000 Vm and f2 = 4000 Vm.
According to Gauss’s theorem,
q
f= (i)
e0
Net flux emerging out of the surface, i.e.,
f = f2 – f1 = 4000 – 2000 = 2000 Vm
By using Eq. (i), we get
q = eof = 8.85 ¥ 10–12 ¥ 2000
= 1.77 ¥ 10−8 C
E xamplE 16 Find the total charge enclosed by a closed surface if number of lines entering is 20,000 and
emerging out is 45000.
Solution Given f1 = 20,000 Vm and f2 = 45,000 Vm.
f = f2 – f1 = net flux emerging out the surface
f = 45000 –20000
= 25,000 Vm
According to Gauss’s theorem
q
f= or q = eof
e0
or q = 8.85 ¥ 10–12 ¥ 25,000
= 22.125 ¥ 10−8 C
E xamplE 17 A point charge of 13.5 Micro Coulomb is enclosed at the centre of the cube of side 6.0 cm. Find
the electric flux (i) through the whole volume and (ii) through one face of the cube.
Solution Given q = 13.5 mC =13.5¥10–6 C and a = 6.0 cm.
(i) According to Gauss’s theorem, the total flux through the whole volume
q
f=
e0
13.5 ¥ 10 -6
=
8.85 ¥ 10 -12
= 1.525 ¥ 106 Nm 2 /C
Since a cube has 6 faces of equal area, the flux through one face of the cube would be
1 q 1.525
= = ¥ 106 Nm 2 /C
6 e0 6
= 2.54 ¥ 105 Nm 2 /C
E xamplE 18 A point charge of 11 Coulomb is located at the centre of a cube of side 5.0 cm. Calculate the
electric flux through each surface.
Solution Given q = 11 C and a = 5.0 cm
As a cube has six faces of equal area, so the flux through each surface of the cube is
1 q 11
= =
6 e 0 6 ¥ 8.85 ¥ 10 -12
= 2.07 ¥ 1011 Nm 2 /C
E xamplE 19 A hollow metallic sphere of radius 0.1 m has 10–8 Coulomb of charge uniformly spread over
it. Determine the electric field intensity (i) on the surface of the sphere (ii) at point 7 cm away
from the centre and (iii) at point 0.5 m away from the centre.
Solution Given radius of the hollow sphere (R) = 0.1 m and charge on it q = 10–8 C.
Formula used for electric field intensity
1 q
E=
4pe 0 r 2
(i) Intensity on the surface of the sphere (r = R) is
1 q
E=
4pe 0 R 2
1 10-8
= ¥
4 ¥ 3.14 ¥ 8.85 ¥ 10-12 (0.1) 2
= 9 ¥ 109 ¥ 10–6
= 9 ¥ 103 N / C
(ii) Intensity at distance 7.0 cm away from the centre. This point lies inside the sphere so that inside the sphere
electric field will be zero, i.e.,
E=0
(iii) Intensity at 0.5m away from the centre
10 -8
E = 9 ¥ 109 ¥
(0.5) 2
E = 0.36 ¥ 103N / C
E xamplE 20 If the charge on a proton is 1.6 ¥ 10–19 Coulomb, find the magnitude of the electric field at a
distance of 1 Å from the proton.
Solution Given qp= 1.6 ¥ 10–19 C and r = 10–10 m.
q 1.6 ¥ 10 -19
E = 9 ¥ 109 ¥ = 9 ¥ 109 ¥
r 2
(10 -10 ) 2
= 1.44 ¥ 1011 V/m
ExamplE 21 Determine the energy gained by an a-particle when it is accelerated through a potential of 1000 volts.
Solution Given qa = 2 ¥ 1.6 ¥ 10–19 C and V = 1000 V.
The energy gained by a-particle is = qV
= 3.2 ¥ 10–19 ¥ 1000 = 3.2 ¥ 10–16 J
E xamplE 22 If the charge on a proton is 1.6 ¥ 10–19 C, find
(i) the electrostatic potential and potential energy at a distance of 1.0 Å from the proton.
(ii) the potential difference between two points 1 Å and 0.2 Å from the proton.
Solution (i) q = 1.6 ¥ 10–19 and r = 1.0 ¥ 10–10 m.
1 q
Potential V=
4pe 0 r
1.6 ¥ 10 -19
= 9 ¥ 109 ¥
10 -10
= 14.4 V
1 q2
Potential energy = -
4pe 0 r
(1.6 ¥ 10 -19 )2
= - 9 ¥ 109
10 -10
-19
= - 23.04 ¥ 10 J
= - 14.4 eV
1 Êq qˆ
(ii) Potential difference = V1 - V2 = -
4pe 0 ÁË r1 r2 ˜¯
È 1 1 ˘
= 9 ¥ 109 ¥ 1.6 ¥ 10-19 Í -10
- -10 ˙
Î 0.2 ¥ 10 10 ˚
= 57.6 V
E xamplE 23 Consider a point charge 15 ¥ 10–6 C. What is the radius of the equipotential surface having
potential 30V?
Solution Given potential = 30 V and q = 1.5 ¥ 10–6 C.
1 q
Potential =
4pe 0 r
1.5 ¥ 10 -6
30 = 9 ¥ 109 ¥
r
1.5 ¥ 10 -16
or r = 9 ¥ 109 ¥ = 450 m
30
E xamplE 24 Calculate the value of poynting vector at the surface of the sun if the power radiated by the sun
is 3.8 ¥ 1026 W and its radius is 7 ¥ 108 m.
Solution Given P = 3.8 ¥ 1026 W and r = 7 ¥ 108 m.
Formula used is
Sav ¥ 4pr2 = P
where Sav is the average poynting vector at surface of the sun
P 3.8 ¥ 1026
Sav = =
4p r 2
4 ¥ 3.14 ¥ (7 ¥ 108 ) 2
= 6.174 ¥ 107 W/m 2
E xamplE 25 Calculate the radiation pressure at the surface of the earth and sun assuming that solar constant
has a value of 2 cal/cm2 min at the surface of the earth and the radius of the sun is 7 ¥ 108 m
and the average distance between earth and sun is 1.5 ¥ 1011 m.
Solution Given
2 cal 2 ¥ 4.2
SE = = -4
cm min 10 ¥ 60
2
J
= 1.4 ¥ 103
m 2 sec
S E 1.4 ¥ 103
[ Prad ] = = = 4.67 ¥ 10-6 N/m 2
c 3 ¥ 108
Further as
Ssr2s = SEr2ES
2
Êr ˆ
2
Ê 1.4 ¥ 1011 ˆ
S S = S E Á ES ˜ = 1.4 ¥ 103 ¥ Á ˜
Ë rs ¯ Ë 7 ¥ 108 ¯
= 5.6 ¥ 107 W / m 2
S s 5.6 ¥ 107
\ [ Prad ]S = =
c 3 ¥ 108
= 0.187 N/m 2
E xamplE 26 Derive Coulomb’s law of electrostatics with the help of Maxwell’s first equation.
Solution From the Maxwell’s first equation
 r   r
div E = or —◊ E = .
e0 e0
Integrating it over a volume V enclosed by a surface S in a medium, we get

  Ê rˆ
Ú —◊ E dV = Ú ÁË e ˜¯ dV
0
V V
But from Gauss’s divergence theorem,
   
Ú —◊ E dV = Ú E ◊ ds
V s
  Ê rˆ 1
\ Ú E ◊ ds = Ú ÁË e 0 ˜¯ dV = e 0 Ú r dV
s V V
  1 q
or Ú E ◊ ds = e 0 ¥ q = e 0
s
If the electric field around the charge is symmetrical, then
q
E ¥ 4p r 2 =
e0
1 q
or E= ¥ 2
4pe 0 r
Hence, the force on a test charge q0 in the electric field E, is
1 qq0
F = q0 E =
4pe 0 r 2
This is the Coulomb’s law.
E xamplE 27 A plane electromagnetic wave propagating along the x-direction has a wavelength 5.0 mm. The
electric field is in the y-direction and its maximum magnitude is 38 V / m. Find the time and
space varying equations for the electric and magnetic fields.
Solution The equations of electric and magnetic fields of a plane electromagnetic wave are given by
E = E0 sin ÈÍ (ct - x)˘˙ H = H0 sin ÈÍ (ct - x)˘˙
2p 2p
and
Îl ˚ Îl ˚
Given E0 =38 V / m and l = 5.0 mm = 5 ¥ 10–3 m.
Hence,
È 2p ˘
E = 38 sin Í -3
(ct - x)˙
Î 5 ¥ 10 ˚
È 2p ¥ 103 ˘
or E = 38 sin Í (ct - x )˙ (i)
Î 5 ˚
The magnitude of the magnetic field is given by
E0 38
H0 = = = 1.27 ¥ 10 -7 Wb/ m 2
c 3 ¥ 108
È 2p ¥ 103 ˘ (ii)
\ H = 1.27 ¥ 10 -7 sin Í (ct - x )˙
Î 5 ˚
The electric field is along y-axis and the magnetic field along z-axis.
E xamplE 28 If the earth receives 2 cal min–1 cm–2 solar energy, what would be the amplitudes of electric and
magnetic fields of radiation.
Solution Here, solar energy which the earth receives is 2 cal min–1 cm–2
E0 m0
\ = = 4p ¥ 4p ¥ 9 ¥ 109 ¥ 10 -7 = 9 ¥ 4p ¥ 4p ¥ 102
H0 e0
= 4p ¥ 3 ¥ 10 = 120p
= 120 ¥ 3.14
= 376.8 ª 377
Poynting vector,
  
P=E¥H
= EH J m -2 sec -1
2 ¥ 4.2
= Jm -2s -1 = 1400 Jm -2 sec -1
60 ¥ 10 -4
\ EH = 1400
E0 E
= 377 =
H0 H
\ E 2 = 1400 ¥ 377
or E = 10 ¥ 14 ¥ 377 = 726.5 V/m

E 726.5
H= = = 1.927
377 377
Amplitudes of electric and magnetic field vectors are
E0 = E 2 = 1024.3 V/m
H 0 = H 2 = 2.717 A/m
 
E xamplE 29 If the magnitude of H in a plane wave is 1 A / m, find the magnitude of E for plane wave in free
space.
H0 e0 Ê m0 ˆ
Solution We know that, = or E0 = Á ˜ H0
E0 m0 Ë e0 ¯
Here H0 = 1 A / m, m0 = 4p ¥ 10–7 Wb / A-m
and e0 = 8.85 ¥ 10–12 C/Nm2
4p ¥ 10 -7
\ E0 = 1 ¥ = 376.72 V/m
8.85 ¥ 10 -12
 ∂r
E xamplE 30 Show that equation of continuity div J + = 0 is contained in Maxwell’s equations.
∂t
Solution According to Maxwell’s fourth equation,

  ∂D
curl H = J +
∂t
Taking divergence of both sides, we get


 Ê ∂D ˆ
div (curl H ) = div Á J + ˜
Ë ∂t ¯

But, div (curl H ) = 0
 
Ê  ∂D ˆ  Ê ∂D ˆ
\ div Á J + ˜ = 0 or div J + div ÁË ˜ =0
Ë ∂t ¯ ∂t ¯
 ∂ 
or div J + (div D) = 0 [\ space and time operations are interchangeable]
∂t
From Maxwell’s first equation,

div D = r,
where r is the surface charge density.

\ div J + =0
t
E xamplE 31 Considering that all the energy from a 1000 W lamp is radiated uniformly, calculate the average
value of the intensity of electric field of radiation at a distance of 2 m from the lamp.
Solution As the total power (P) is radiated uniformly, energy flux per unit area per second at a distance r from the lamp
is given by
P0 1000
S= 2
= W/m 2
4p r 4p (2) 2
  
From the poynting theorem, | S | = | E ¥ H | = EH sin 90∞
 E and H are perpendiculars to each other,
1000
\ EH = (i)
16p
E m0
But = = 376.72 ohm (ii)
H e0
Multiplying Eqs. (i) and (ii), we get
1/ 2
E 376.72 ¥ 1000 Ê 376.72 ¥ 1000 ˆ
EH ¥ = fi E=Á = 86.59 V /m
H 16p Ë 16 ¥ 3.14 ¯˜
E xamplE 32 If the relative permittivity of distilled water be 81. Calculate refractive index and velocity of
light in it.
me c
Solution We know that, m ¢ = and v =
m 0e 0 m¢
Here e = 81e0 and for distilled water m ª m0
3 ¥ 108
\ m ¢ = 81 = 9 and v = = 3.33 ¥ 107 m/sec
9
E xamplE 33 Consider an infinite conducting sheet in the xy-plane with a time dependent current density kti,
m k
where k is constant. The vector potential at (x, y, z) is given by A = 0 (ct – z)2. Find the magnetic field B.
4c
Solution  
B=—¥ A
iˆ ˆj kˆ
∂ ∂ ∂
=
∂x ∂y ∂z
m0k
(ct - z ) 2 0 0
4c
m k ∂
= 0 (ct - z ) 2 ˆj
4c ∂ z
0k
(z ct ) ˆj
2c

E xamplE 34 A magnetic field B = B0 (iˆ + 2 ˆj - 4kˆ) exists at a point. If a test charge moving with a velocity

v = v0 (3iˆ - ˆj + 2kˆ) experiences no force at a certain point, what will be the electric field at that point in SI
units?
Solution Lorentz force is given by
   
F = qE + q (v ¥ B )
= 0 (as per question)
  
\ E = - (v ¥ B )
iˆ ˆj kˆ

E = v0 B0 1 2 -4
3 -1 2

E = - v0 B0 [iˆ(0) - ˆj (14) + kˆ (-7)]

E 7v0 B0 (2 ˆj kˆ )V /m
 
E xamplE 35 Find the electric and magnetic fields E ( z , t ) and B( z , t ) , respectively, corresponding to the
scalar potential f(z, t) = 0 and vector potential A = zt iˆ
Solution
iˆ ˆj kˆ
   ∂ ∂ ∂
B=—¥ A=
∂x ∂y ∂z
tz 0 0
 ∂
B = (tz ) = t
∂z
 ∂A ∂
E =f - =0- ( zt ) = - ziˆ
∂t ∂t

E ziˆ V /m
E xamplE 36 A plane polarized electromagnetic wave in free space at time t = 0 is given by


E (x, z) = 10 ˆj exp [i(6x + 8y)]. Find the corresponding magnetic field.
 
Solution From expression of E , k = 6iˆ + 8 ˆj
 
  ∂B   ∂H
Maxwell’s equation — ¥ E = - or — ¥ E = - m0 gives
∂t ∂t
    k ˆ 
ik ¥ E = m0iw H fi H = k¥E
m0w
 1 ˆ 
or H= k¥E ( w/k = c in free space)
m 0c
10 È (6iˆ + 8 ˆj ) ˘
= Í ¥ ˆj ˙ exp[i (6 x + 8 y )]
m0c Î 10 ˚
6
exp( i (6 x 8 y ))kˆ A/m
0c
E xamplE 37 A plane electromagnetic wave has magnetic field

 È k ˘
B( x, y, z , t ) = B0 sin Í( x + y ) + w t ˙ kˆ
Î 2 ˚
where k is wave number and iˆ, ˆj and kˆ are cartesian unit vectors in x, y and z directions respectively. Find

the electric field E (x, y, z, t) and average poynting vector.
  
Solution Using k ¥ H = - e 0w E
 1  
E=- k¥H
e 0w
B0 k È k ˘ È (iˆ + ˆj ) ˆ ˘
=- sin Í( x + y ) + wt ˙ Í ¥ k˙
m0e 0w Î 2 ˚Î 2 ˚
B0 k È k ˘ Ê - ˆj + iˆ ˆ
=- sin Í( x + y ) + wt ˙ Á ˜
m0e 0w Î 2 ˚Ë 2 ¯
B0 k k iˆ ˆj
E sin ( x y) t units
0 0 2 2
Average poynting vector is given by

   B
<P > = E ¥ H = E ¥
m0
- B02 k È k ˘ Ê iˆ - ˆj ˆ ˆ
= sin 2 Í( x + y ) + wt ˙ Á ˜ ¥k
˚Ë 2 ¯
2
m0 e 0w Î 2
- B02 k Ê - ˆj - iˆ ˆ
<P > = Á ˜
2 m02e 0w Ë 2 ¯
B02c 2 k Ê iˆ ˆj ˆ
= Á ˜ units
2 0 Ë 2 ¯
E xamplE 38 In a non-conducting medium characterized by e = e0, m = m0 and conductivity s = 0, the electric


field (in Vm ) is given by F = 20 sin[108t – kz] ĵ . Find the magnetic field, H (in Am–1).
–1
Solution For non-conducting medium,

 1  
H= (k ¥ E )
mw
 20k
H= (kˆ ¥ ˆj )sin[108 t - kz ]
m108
20k
H sin[108 t kz ]iˆ A/m
108
E xamplE 39 Consider two concentric conducting spherical shells with inner and outer radii a, b, c and d, as
shown. Both the shells are given Q amount of positive charges. Find electric field in different regions.
Solution Conducting shells are given, means the supplied charge Q will be distributed over their outer surfaces (see the
figure). Moreover, a charge — Q will be induced on the inner surface of the outer shell due to the positive charge Q on
the inner shell. Due to induction, a charge +Q will be induced on the outer surface of the outer shell. It means the total
charge on its outer surface will be as 2Q.
Field in region I (r < a) + + + + Q+ + Q = 2Q
There is no charge enclosed by the Gaussian surface within this region. ++ +
+
  q ––
+
V
Hence, Gauss’s law reads  Ú E ◊ dS = eenc0 = 0
++++++
++
+ Q+ IV
+ III
 II
+++++
fi
+ +
E = 0 for r < a
+ + +
I
a
Field in region II (b < r < a) d b
+
In this region also, the charge enclosed by the Gaussian surface would be zero. + +
c
++

fi E = 0 for a < r < b.
+
+ + –Q
Field in region III (b < r < c)
++
++ ++++
Here enclosed charged will be as Q.
 
\ Gauss’s law reads  Ú E ◊ dS = Q /e 0
 Q
or E4pr2 = Q/e0 fi E = rˆ
4pe 0 r 2
Field in region IV (c < r < d)
In this region, the enclosed charge = –Q + Q = 0

fi The field E = 0.
Field in region V (r > d)
The charge enclosed by Gaussian surface in this region will be 2Q – Q + Q = 2Q
  2Q
\ Gauss’s law reads  Ú E ◊ dS = e 0
 2Q Q
fi E= rˆ = rˆ
4pe 0 r 2 2pe 0 r 2
E xamplE 40 Consider two concentric uniformly charged spherical shells with inner and outer radii a, b,
c and d, as shown. Both the shells carry equal amount of positive charge Q. Find electric field in different
regions.
Solution Charge density of inner spherical shell IV

Q
r1 =
4
p (b3 - a 3 ) b
3 c
Q
Charge density of outer spherical shell r2 =
4 a q
p (d 3 - c3 )
3 d I
Region I (r < a)
Charge enclosed by Gaussian surface for r < a will be zero. II
   III
Ú E ◊ dS = 0 fi E = 0 .
Hence, Gauss’s law 
IV
Region II (a < r < b) V
4
In this region, the charge enclosed by Gaussian surface will be r1 ◊ p (r3 – a3)
3
Hence, Gauss’s Law reads
  Q 4
Ú E ◊ dS = 4 3 3 ◊ 3 p (r - a )
3 3
e 0 p (b - a )
3
Q ( r 3 - a3 )
fi E ◊ 4p r 2 =
e 0 (b3 - a 3 )
or
 Q(r 3 - a3 )
E= rˆ
4pe 0 r 2 (b3 - a 3 )
Region III (b < r < c)
Charge enclosed by Gaussian surface will be Q.
 Q
fi E= rˆ
4pe 0 r 2
Region IV (c < r < d)
Charge enclosed by Gaussian surface will be
4
r2 ◊ p (r3 – c3)
3
\ Gauss’s law reads
  Q 4
Ú E ◊ dS = 4 3 3 ◊ 3 p (r - c )
3 3
e 0 p (d - c )
3
Q ( r 3
- c3 )
or E 4p r 2 =
e 0 (d 3 - c3 )
 Q(r 3 - c3 )
or E= rˆ
4pe 0 r 2 (d 3 - c3 )
Region V (r > d)
Charge enclosed by Gaussian surface will be
Q + Q = 2Q
  2Q
Hence, Ú E ◊ dS = e 0
 2Q
or E= rˆ
4pe 0 r 2
Q
= rˆ
2pe 0 r 2
E xamplE 41 The electric field of a uniform plane wave propagating in a dielectric, non-conducting medium

is given by E = xˆ 10 cos(6p ¥ 107t – 0.4pz) V/m. Find phase velocity of wave.
 
Solution For the plane wave whose electric field is given by E = E0 cos (wt – kz), the phase velocity is vp = w/k. Hence
w 6p ¥ 107
vp = = = 1.5 ¥ 108 m/s
k 0.4 ¥ 3.14
E xamplE 42 Determine penetration depth by which an electromagnetic wave enters into copper,
if rcu = 1.69 ¥ 10–8 mm and frequency = 104 MHz.
Solution
2 r
d= =
msw pfm
1.69 ¥ 10-8 ¥ 10-6
=
3.14 ¥ 1010 ¥ 4p ¥ 10-7
= 0.0654 ¥ 10–8 m = 0.654 nm
E xamplE 43 In free space, H = 0.1 cos (2 ¥ 108t – kx) ĵ H/m. Calculate k, l and T.
Solution
w
k=
v
For free space, v = c
w 2 ¥ 108
k= = = 0.667 rad/m
c 3 ¥ 108
2p 2p 2 ¥ 3.14
k= ,l = = = 9.42 m
l k 0.667
2p 2 ¥ 3.14
T is the period, T = = = 3.14 10 8 s
w 2 ¥ 108

E xamplE 44 In a lossless medium for which Z = 60p, mr = 1, and H = –0.1 cos (wt – z) iˆ + 0.5 sin (wt – z) A/m.
Calculate er, w and E.
Solution
m m0 m r 120p
Z= = =
e e0 er er
120p
fi er = = 2 fi er = 4
60p
w 2w
b = w me = w m0e 0 m r e r = 4=
c c
b c 3 ¥ 108
w= = = 1.5 ¥ 108 rad/s (  b = 1 due to b z = z )
2 2

  ∂E  1  
—¥H =e fi E = Ú — ¥ H dt
∂t e
iˆ ˆj kˆ
  ∂H y ∂H x ˆ
— ¥ H = ∂ /∂x ∂ /∂y ∂ /∂z = - iˆ + ¥ j
∂z ∂z
Hx Hy 0
= 0.5 cos(w t - z )iˆ - 0.1 sin(w t - z ) ˆj
 1
Hence E = Ú [+ 0.5 cos(w t - z )idt ˆ - 0.1 sin(w t - z ) ˆjdt ]
e
0.5 0.1
= sin(w t - z )iˆ + cos(w t - z ) ˆj
ew ew
 0.5 0.1
or E= sin(w t - z )iˆ + cos(w t - z ) ˆj
e 0e r w e 0e r w
0.5 0.1
= sin(w t - z )iˆ + cos(w t - z ) ˆj
8.85 ¥ 10-12 ¥ 4 ¥ 1.5 ¥ 108 8.85 ¥ 10-12 ¥ 4 ¥ 1.5 ¥ 108

fi E 94.16 sin( t z )iˆ 18.83 cos( t z ) ˆj
E xamplE 45 If er = 1, mr = 20 and s = 3 mhos/m for a medium and the electric field of an electromagnetic
wave is

E = 2e a z sin(108 t - b z ) ˆj V/m

Find a, b and H .
s 3
Solution As =
ew 8.85 ¥ 10-12 ¥ 108
= 3389 >> 1.
the medium is good conductor at frequency of operation.
mws
fi a =b=
2
1/2
È 4p ¥ 10-7 ¥ 20 ¥ 108 ¥ 3 ˘
=Í ˙
Î 2 ˚
b = 61.4 rad/m
mw 4p ¥ 10-7 ¥ 20 ¥ 108
h= =
s 3
800p
= W
3
Ê pˆ
H = H 0e -a z sin Á w t - b z - ˜
Ë 4¯
E0 3
H0 = =2 = 69.1 ¥ 10-3 A/m
h 800p
 Ê pˆ 
H = 69.1 ¥ e -61.4 z sin Á108 t - 61.4 z - ˜ (k ¥ ˆj )
Ë 4¯
 Ê pˆ
H = - 69.1e -61.4 z sin Á108 t - 61.4 z - ˜ iˆ mA/m
Ë 4¯
E xamplE 46 In non-magnetic medium


E = 4 sin (2p ¥ 107t – 0.8x) k̂ V/m
Find power crossing 100 cm2 of plane
2x + y = 5
Find impedance (h), and er.
Solution
w er
b = w me = w m0e 0e r =
c
b c 0.8 ¥ 3 ¥ 10 12
8
er = = =
w 2p ¥ 107 p
er = 14.59
m m0 120p 120p ◊ p
h= = = =
e e 0e r er 12
= 10p 2 W
   E2
P = E ¥ H = 0 sin 2 (w t - b x)iˆ
h
T
1  E2
Pav = Ú
tn
Pdt = 0 iˆ
2h
16
= iˆ
2 ¥ 10p 2
2iˆ + ˆj
On plane 2x + y = 5, the normal is nˆ =
5
  
\ Pavg = Ú Pav ◊ dS = Pa ◊ san
Ê 100 ¥ 10-4 ˆ
= (81 ¥ 10-3 aˆ x ) ◊ Á (2iˆ + ˆj )˜
Ë 5 ¯
= 724.5 W
E xamplE 47 Consider two concentric spherical conducting shells (negligible thickness) having equal charge
q. The radius of inner shell is “a” and that of the outer shell is “b”. Calculate the electric field in three regions
(a) r < a
(b) a < r < b
(c) r > b
Solution In Figure, dark black circles represent the spherical shells, which carry equal
charge q, and grey color circles represent the Gaussian surfaces in all the three regions.
(i) For interior region (r < a)
a
Charge enclosed by the Gaussian surface drawn in this region is zero. Using
Gauss’s law for spherical symmetry b
  q
Ú E ∑ dA = eenc0
Since the field E would be the same on the Gaussian surface

0
E ◊ 4p r 2 =
e0
E=0
(ii) For middle region (a < r < b)
Charge enclosed by the Gaussian surface drawn in this region is q. Using Gauss’s law, we get
q
E ◊ 4p r 2 =
e0
 q
E= rˆ
4pe 0 r 2
(iii) For outer region (r > b)

Charge enclosed by the Gaussian surface drawn in this region is 2q. Using Gauss’s law, we get
  q
Ú E ∑ dA = eenc0
2q
E ◊ 4p r 2 =
e0
 q
E= rˆ
2pe 0 r 2


Q.1 The direction of —f is always
(a) parallel to the surface (b) perpendicular to the surface
(c) depends upon the shape of surface (d) none of these

Q.2 The vector field A is solenoidal if
    
(a) — ¥ A = 0 (b) A = 0 (c) — ◊ A = 0 (d) none of these

Q.3 A vector field A is conservative if
     
(a) — ¥ A π 0 (b) — ◊ A = 0 (c) A = —f (d) none of these
 
Q.4 The vector function — ◊ A represents
(a) the total flux over any arbitrary closed surface
(b) the inward flux density at the point (x,y,z)
(c) the outward flux density at the point (x,y,z)
(d) none of these
Q.5 The vector field is irrotational if
       
(a) — ¥ A = 0 (b) — ◊ A = 1 (c) — ◊ A = 0 (d) — ¥ A = 1
   
Q.6 If a is a constant vector then — ¥ (a ¥ r ) is equal to
  
(a) zero (b) 2a (c) a /2 (d) a

Q.7 For a conservative field E
       
(a)  Ú E ◊ dl = 0 (b) Ú ◊ ds = 0
E (c) — ◊ E π 0 (d) — ¥ E π 0
Q.8 The divergence of curl of a vector is always

(a) 1 (b) p/2 (c) 1/2 (d) zero
Q.9 The unit of e0 in SI system of units is
(a) Nm2 / C2 (b) C / S (c) C2 / Nm2 (d) N / C
Q.10 The unit of electric flux in SI system of units is
(a) Weber (b) Gauss (c) Nm2 / C (d) N / C
Q.11 A charge of 15 mC is placed at one corner of a cube. The electric flux will be
(a) zero through each face
(b) same through each face
(c) zero through three faces meeting at the location
(d) none of these.
Q.12 The variation of electric field intensity E with distance r from the centre of hollow spherical shell of
radius R is
E E E E
r r r r
(a) (b) (c) (d)
Q.13 The electric field intensity E inside a uniformly charged sphere varies with distance r of the observation
point as
1 1
(a) E µ r (b) E µ (c) E µ r2 (d) E µ 2
r r
Q.14 The electric field between two oppositely charged plates having equal charge density s is given by
(a) s / e0 (b) s /2e0 (c) zero (d) 2s /e0
Q.15 Which of the following is zero
(a) grad div (b) div grad (c) curl grad (d) curl curl
Q.16 The relation between electric field and potential is
   
(a) — ◊ E = V (b) E = –—2 V (c) E = –— V (d) E = —2 V
Q.17 The work done in displacing a charge 2C through 0.5 m on an equipotential surface is
(a) zero (b) 4 J (c) 1 J (d) none of these
Q.18 A point charge q is located at the origin. The amount of work done in bringing a unit positive charge
from infinity to the origin is
(a) zero (b) finite (c) infinite (d) none of these
Q.19 Which of the following equations tells us about the non-existence of the magnetic monopole?
 
 ∂B     ∂B
(a) curl E = - (b) div B = 0 (c) div D= r (d) curl H = J +
∂t ∂t
Q.20 Displacement current appears because of
(a) time varying electric field (b) time varying magnetic field
(c) negative charge only (d) positive charge only
Q.21 Skin depth is the distance

(a) travelled by EM wave in conductor
(b) by which electric field amplitude of EM wave falls to 1e times of its initial value
(c) by which wavelength of EM wave falls to 1e times of its initial value
(d) by which wavelength of EM wave falls to 12 times of its initial value
Q.22 EM wave attenuation is
(a) sudden increase in wave amplitude (b) very slow increase in wave velocity
(c) exponential decay of wave amplitude (d) exponential increase of wave velocity
Q.23 Dispersion relation of an EM wave is
(a) v = fl (b) v = w/k
(c) f = w/2p (d) relation between w and k
Q.24 For wave attenuation
(a) real part of wave number is responsible (b) imaginary part of wave number is responsible
(c) real part of frequency is responsible (d) imaginary part of frequency is responsible
Q.25 The wave equation for EM wave propagation in a conductor is
  
(a) relation between E and B (b) relation between E and w
 
(c) time variation of E and B fields (d) a relation showing time and space variation of either
 
E or B
Q.26 Stationary current is the current
(a) which does not flow
(b) produced by holes
(c) where charge density remains constant during motion of charges
(d) where charge density vanishes
Q.27 In EM wave
(a) electrons produce magnetic field only
(b) electron produce electric field only
(c) time variation of electric field produces magnetic field and vice-versa
(d) time variation of electric field guides the wave
Q.28 Following expression represents the wave motion
(a) E = E0 sin wt (b) E = E0 sin (wt – kx)
(c) E= E0 cos kx (d) E = E0 sin cos wt
Q.1 Discuss conservative versus nonconservative fields by giving examples.

Q.2 How do you find the magnetic flux crossing an infinitesimal surface?
Q.3 Provide physical interpretations for the closed surface integrals of any two vectors.
Q.4 What are different ways in which an emf is induced around a loop?
Q.5 State Ampere’s circuital law.
Q.6 What is displacement current?
Q.7 State Gauss’s law for electric field.
Q.8 State Gauss’s law for magnetic field.
Q.9 What is the physical interpretation of Gauss’ law for magnetic field?
Q.10 State the law of conservation of charge.
Q.11 How is Gauss’s law dependent on Ampere’s circuital law?
Q.12 What is Gaussian surface?
Q.13 What is the electric flux through a closed surface surrounding a dipole?
Q.14 Can we apply Gauss’s law to calculate the electric field due to electric dipole? Explain.
Q.15 A Gaussian surface encloses no net charge. Does it mean E = 0 on its surface.
Q.16 What do you understand by pointing vector?
Q.17 Write Maxwell’s equations.
Q.18 State Faraday’s laws of electromagnetic induction.
Q.19 What do you mean by a waveguide?
Q.20 Write an expression for characteristic impedance of a co-axial cable.
P ractice P roblems
general questions
Q.1 Describe gradient of a scalar field in Cartesian coordinates. Explain its physical significance.
Q.2 The gradient of a scalar field is a vector. Hence explain how can you produce a vector from a scalar
field.
Q.3 Give the physical interpretation of grad V.
Q.4 Define divergence of a vector field. What is its physical meaning? Give two examples.
Q.5 Divergence of a vector field is a scalar quantity. Hence explain how you can produce a scalar field from
a vector field.
Q.6 Derive an expression for divergence of a vector field in Cartesian coordinates from first principle.
 
Q.7 What do you mean by a solenoidal vector field? Give one example. What is the meaning of — ◊ E π 0?
Q.8 State and prove Gauss’s divergence theorem.

Q.9 Prove that the volume integral of the divergence of a vector field A taken over any volume is equal to

the surface integral of A over the closed surface surrounding the volume.
Q.10 Define curl of a vector field and give its physical significance. Show that curl of a vector field is a
vector quantity.
Q.11 Calculate the value of the curl of a vector in terms of Cartesian coordinates.
Q.12 What is an irrotational field? Give one example.
Q.13 Prove that the curl of linear velocity of the particles of a rigid body rotating about an axis passing
through it is twice the angular velocity.
    1  
Q.14 If w ¥ r = V , prove that w = curl V , where w is a constant vector.
2
Q.15 Show by actual computation that curl gradient of a scalar function is always zero or curl of grad f = 0.
Q.16 Show that the curl of a uniform electric field is zero.
Q.17 Show that a vector field whose curl is everywhere zero can be expressed as the gradient of another
suitable scalar field. What is this type of field called?

Q.18 Prove that div curl A = 0.
 
Q.19 If a vector B is curl of another vector A, then prove that the divergence of such vector is zero.
Q.20 Show that a vector field whose divergence is everywhere zero can be expressed as curl of some other
suitable vector field.
Q.21 State and prove Stokes’ theorem. Discuss its importance.
Q.22 What is a conservative field? Show that a conservative field is the gradient of a scalar field and curl of
such a field is zero.

Q.23 Show that electric field is conservative and curl E = 0.
Q.24 What is the difference between a conservative and non-conservative field? Give one example of each.
Q.25 What do you understand by the term charge density?
Q.26 What is line charge density? Derive an expression for the electric field due to an infinitely long
uniformly charged straight wire using Coulomb’s law.
Q.27 A thin non-conducting rod of length l carries a positive charge distributed uniformly over its length. If
the linear charge density is l, find the intensity of the electric field at a point at a distance a from the
near end of the rod and on its axis.
Q.28 Two parallel infinite wires have uniform line charge densities l1 and l2 separated by a distance x.
Calculate the electric force per unit length on one wire as a result of the other.
Q.29 Derive an expression for electric field at a point situated on the axis of a uniformly charged ring.
Q.30 Define surface charge density and volume charge density. State the relation between electric intensity
and charge density.
Q.31 Is volume charge density invariant (under Lorentz transformations)?
Q.32 Find the electric field due to a circular charged disc at a point on a line perpendicular to the disc and
passing through its centre. Hence calculate electric field due to an infinitely large plane conducting
sheet of charge.
Q.33 Calculate the electric field strength due to a uniform charged circular sheet on the axis.
Q.34 Explain the meaning of the term electric flux. What are its dimensions and S.I. units.
Q.35 State and prove Gauss’s theorem in electrostatics. Prove that total flux over a surface due to a charge
lying outside is zero.
Q.36 State Gauss’s theorem. Derive the differential form of this Gauss’s theorem.
Q.37 Write the law for a volume distribution of charge.
Q.38 Apply Gauss’s theorem to calculate the electric field due to a uniformly charged solid cylinder.
Q.39 Prove that the electric field at a point inside a uniformly charged cylinder of infinite length is
proportional to the distance of the point from the axis.

Q.40 Apply Gauss’s theorem to find the electric field strength E near a plane non-conducting thin sheet of
charge of infinite extent. Hence show that the field is independent of the distance of the observation
point from the sheet.
Q.41 Apply Gauss’s law to calculate
(i) The electric field at any point due to two parallel sheets of charge.
(ii) Calculate the intensity of the electric field at a point between oppositely charged parallel plates.
Q.42 Using Gauss’s theorem calculate the electric field due to a uniform spherical shell of charge at a point
(i) Outside the shell and (ii) inside the shell. Hence show that for points lying external to it a uniformly
charged spherical shell behaves as if the entire charge were concentrated at its centre and for point
lying inside it the electric field is zero.
Q.43 Using Gauss’s theorem calculate the electric field due to a uniformly charged non-conducting solid
sphere at a point
(i) Outside the sphere
(ii) On the surface of the sphere, and
(iii) Inside the sphere.
Q.44 State
  and prove Gauss’s law or Gauss’s theorem. Express it in differential form and show that
—◊ E = r/ e 0
Q.45 Show that Coulomb’s law can be deduced from Gauss’s law and considerations of symmetry.
Q.46 Coulomb’s law is a special case of Gauss’s law. Explain.
s
Q.47 Prove that the electric field on the surface of a conductor is where s is the surface charge density.
e0
Hence find the electric field near a charged conducting sheet with the same surface density of charge.
Q.48 State and
 prove Ampere’s circuital law of magnetic field. Deduce Ampere’s law in the form
Ú B ◊ dl = m0 I , where the symbols have their usual meaning.
Q.49 Show that the line integral of the magnetic field over a closed path is independent of the shape of the
path.
Q.50 Using Ampere’s law obtain an expression for the magnetic field due to a current carrying straight
conductor of infinite length.
Q.51 Using Ampere’s law calculate the magnetic field at point inside a long current carrying solenoid.
Q.52 Explain the concept of Maxwell’s displacement current and show how it led to the modification of
Ampere’s law.
Q.53 What is equation of continuity? Explain it. How could Maxwell correct and present Ampere’s law in
its generalized form?
Q.54 What are Maxwell’s equations? Derive Maxwell’s equations (differential form). Discuss integral form
of above equations. What are the significance of these equations to electricity and magnetism?
Q.55 Obtain the electromagnetic wave equations, using Maxwell’s equation, in an isotropic dielectric
medium and show that the speed of wave is less than its speed in vacuum.
Q.56 Obtain Maxwell’s equations and deduce an expression for the velocity of propagation of a plane
electromagnetic wave in a medium of dielectric constant e and relative permeability m.
Q.57 Define poynting vector. Derive an expression for it and explain its physical significance for
electromagnetic wave in free space.
Q.58 Derive the electromagnetic wave equation from Maxwell field equations. Consider plane wave solutions
of this equation and prove that the energy density associated with such a wave in a stationary homogeneous
non-conducting medium propagates with the same speed with which the field vectors do.
Q.59 A plane monochromatic electromagnetic wave propagates in a conducting medium. Show that
attenuation is equal to phase vector.
Q.60 Discuss the propagation of plane monochromatic electromagnetic waves in conducting media. Derive
the dispersion equation and thus obtain: (i) phase velocity (ii) refractive index (iii) skin depth.
Q.61 Show that inside the conducting medium the wave is damped and obtain an expression for the skin
depth d.
Q.62 What is a wave guide? Describe the propagation of electromagnetic wave along a hollow wave guide
of uniform cross section.
Q.63 Give a brief note on coaxial cables with special reference to characteristic impedance.
Q.64 Write note on
(i) Displacement current
(ii) Poynting vector
Q.65 What do you understand by dispersion relation of an electromagnetic wave?
Q.66 Using dispersion relation, find the wavelength and phase velocity of an electromagnetic wave in a
dielectric medium having e = 4e0 and m = m0.
Q.67 What do you understand by wave attenuation? Find an expression for the skin depth of an electro-
magnetic wave in a conductor.
Q.68 Discuss in detail the electromagnetic wave propagation in a conducting medium. What are the roles of
real and imaginary parts of the wave number of the wave?
Q.69 Find the time taken by a charge Q placed in the interior of a copper to drop to 36.8 percent of its initial
value. Take s = 5.8 ¥ 107 mhos/m and e = e0.
Q.70 Plot a graph between the skin depth in copper and frequency of an electromagnetic wave when
s = 5.8 ¥ 107 mhos/m, and m = m0.
Q.71 Prove that electric and magnetic field vectors of an electromagnetic wave do not remain in phase when

it propagates in a conductor. Also show that qH – qE = qK, where qH is the phase of magnetic field H ,
 
qE is the phase of electric field E and qK is the phase of wave vector k .
Theory of Relativity 12
Learning Objectives
L01 Understand the inertial and L04 Explain length contraction, time
non-inertial frames of reference dilation
L02 Learn about Galilean transformation L05 Discuss addition of velocities
and Michelson-Morley experiment LO6 Evaluate variation of mass with velocity
L03 Know postulates of special theory of and Einstein’s mass energy relation
relativity and Lorentz transformation
Introduction
Before the beginning of the 20th century, main branches of physics, namely mechanics and electromagnetism,
had developed independently. It was a firm opinion of the physicists that these two have no strong relation
with each other. However, early in the 20th century they started facing many new and basic problems. For
example, Newton’s second law of motion did not give correct results when applied to the objects moving
with high speed comparable to the speed of light. Moreover, they noticed that for two observers, which
are in relative motion, the same set of transformation equations cannot be used to transform the laws of
mechanics and electromagnetism from the frame of reference of one observer to the frame of reference
of the other observer. However, the introduction of special theory of relativity by Einstein in 1905 led to
the solution of these and other difficulties. The special theory of relativity deals with the objects or frames
of references, which are moving with uniform velocity relative to each other. For dealing with accelerated
frames of references, Einstein developed the general theory of relativity in 1915.
12.1 Frame OF reFerence LO1

It is well known that the motion of a body is better understood if it is described with respect to some well
defined coordinate system. This coordinate system is called frame of reference. Thus, a system of coordinate
axes which defines the position of a particle in two-or three-dimensional space is called frame of reference.
It is divided into two categories:

(i) Inertial i.e., unaccelerated system
(ii) Non-inertial i.e., accelerated system
12.1.1 Inertial Frame of reference
Newton’s first law of motion says that a system at rest will remain at rest or a system in uniform motion will
remain in uniform motion till an external force is applied on it. The systems in which the law of inertia holds
good are called inertial systems or inertial frames of reference. Actually the earth is not an inertial frame of
reference because of its orbital and rotational motions. However, for most of the purposes, the earth may be
regarded as an inertial frame of reference.
12.1.2 non-inertial Frame of reference

A non-inertial system is the one in which Newton’s first law of motion does not hold good. A frame of reference
which is in accelerated motion with respect to an inertial frame is called non-inertial frame of reference.
12.2 GaLILean TransFOrmaTIOn LO2

There are the transformations which can transform the Y F Y¢ F¢ P(x, y, z, t)
coordinates of a particle from one inertial system to another. (x ¢, y ¢, z ¢, t ¢)
x
Consider two inertial systems F and F¢, where F¢ is moving
x¢
with uniform velocity v relative to F along the +ve direction of
x-axis (Fig. 12.1). We further consider that the origin of both v
the systems coincide at time t = t¢ = 0. Let an event happen at
O O¢ X, X ¢
point P whose coordinates are (x, y, z, t) and (x¢, y¢, z¢, t¢) with
respect to the frames of references F, and F¢ respectively. It
can easily be seen that these coordinates are related, as Z Z¢
Figure 12.1
x ¢ = x - vt ¸
y¢ = y Ô
Ô (i)
˝
z¢ = z Ô
t¢ = t Ô˛
Eq. (i) is known as Galilean transformation for position.

Now, consider that frame F¢ is moving along any direction with velocity v relative to F such that
    Y Y¢ P(x, y, z, t)
v = vx iˆ + v y ˆj + vz kˆ (x ¢, y ¢, z ¢, t ¢)
vxt F¢
where vx, vy and vz are the components of v along X, Y and
O´
Z-axes, respectively, as shown in Fig. 12.2. Suppose the origins F
X´
vyt
of the two systems F and F¢ coincide at t = t¢ = 0. Let (x, y, z,
t) and (x¢, y¢, z¢, t¢) are the coordinates of the event happening O X
at point P. At the time of event the frame F¢ is separated from
vzt
frame F by a distance vxt, vyt and vzt along X, Y and Z-axes Z¢
Z
respectively. Then we have,
Figure 12.2
x ¢ = x - vx t ¸
y ¢ = y - v y t ÔÔ (ii)
˝
z ¢ = z - vz t Ô
t¢ = t Ô˛
Eq. (ii) is also known as Galilean transformation for position.

Galilean velocity transformation of the particle can be obtained by differentiating Eq. (ii), with respect to
d d
time. By using = and vx, vy and vz to be constant, we obtain
dt dt ¢
dx ¢ dx ¸
= - vx Ô
dt ¢ dt
Ô
dy ¢ dy Ô (iii)
= - vy ˝
dt ¢ dt Ô
dz ¢ dz Ô
= - vz Ô
dt ¢ dt ˛
u x¢ = u x - vx ¸
or u ¢y = u y - v y Ô˝ (iv)
u z¢ = u z - vz Ô˛
where ux, uy and uz are the velocities of the particle observed by an observer O in system F and u¢x, u¢y and
u¢z are the velocities of the particle observed by O¢ in system F¢ along X, Y and Z-axes, respectively. From
Eq. (iv), we have,
u x¢ iˆ + u y¢ ˆj + u z¢ kˆ = u x iˆ + u y ˆj + u z kˆ - vx iˆ - v y ˆj - vz kˆ
  
or u¢ = u - v (v)
where iˆ, ˆj , kˆ are unit vectors along X, Y and Z-axes, respectively. Eq. (v) represents the Galilean transformation
of velocity of particle.
Similarly, Galilean acceleration transformation of the particle can be represented by the following equations
by knowing the fact that the acceleration of a particle is the time derivative of its velocity.
du x du y du z
ax = , ay = and az =
dt dt dt
To find the Galilean acceleration transformations, we differentiate the velocity transformation and use the
fact that t¢ = t and vx, vy and vz are the constants. This yields
ax¢ = ax, ay¢ = ay and az¢ = az

12.3 Michelson-Morley experiMent LO2

In 19th century, scientists had assumed that a hypothetical medium
M1
called luminiferrous ether is required for the propagation of the
light. It was considered that the ether exists uniformly in the
space and it is at rest relative to the earth and other planets. The
basic purpose of Michelson-Morley experiment was to confirm
l Beam II
the existence of stationary ether. The existence of stationary
ether (an absolute frame of reference) can be confirmed if we P P´
can measure the absolute velocity of earth with respect to the
stationary ether. S
M2
Figure 12.3 shows an arrangement of Michelson-Morley
experiment. Its main components are mirror M1 and M2,
l
glass plates P and P¢ of the same size, a source of light S and
a telescope T. The light beam moves to the mirror M1 and the Beam I
other moves to mirror M2. These two beams are reflected back T
by these mirrors and again are recombined at plate P. Finally,
they enter the telescope and produce interference. In the context Figure 12.3
of interference, this experiment has been discussed in detail in
Chapter 1.
In this experiment, the mirrors M1 and M2 are set such that PM1 = PM2 = l. According to Galilean
transformation, the velocity of light in a frame moving with constant velocity v relative to stationary either
from P to M2 is (c – v) while from M2 to P is (c + v). If t1 be the time taken by the transmitted beam from P
to M2 and back, then
l l
t1 = + (i)
c-v c+v
2lc 2l È 1 ˘
or t1 = 2 = Í 2 2˙
c -v 2
c Î1 - v /c ˚
-1
2l È v 2 ˘
or t1 = Í1 - 2 ˙
c Î c ˚
By using Binomial expansion theorem for the case of v/c << 1, we get
2l È v 2 ˘
t1 = Í1 + ˙ [Neglecting higher order terms] (ii)
c Î c2 ˚
The distance travelled by the light beam in time t1 is given by
È v2 ˘
x1 = t1c = 2l Í1 + 2 ˙ (iii)
Î c ˚
If t2¢ be the time taken by beam-II from P to M1 and in the same time distance travelled by this beam is ct2¢.
In this time t2¢, the mirror M1 shifts to M1¢ and travels a horizontal distance vt2¢.
With the help of Fig. 12.4, we get PO = l, PO¢ = ct2¢ and OO¢ = vt2¢. M1 M´1
O O´
( PO ¢ ) 2 = ( PO ) 2 + (OO ¢ ) 2
(ct2¢ ) 2 = (l ) 2 + (vt2¢ ) 2
or (c 2 - v 2 )t2¢ 2 = l 2
1/2
Ê l2 ˆ
t2¢ = Á 2 ˜
Ë c - v2 ¯
P
A P´
1 1
or t2¢ = ¥ (iv)
c 1 - v2 / c2 Figure 12.4
2l 1
Total time, t2 = 2t2¢ = ¥
c 1 - v 2 /c 2
-1/2
2l Ê v2 ˆ
or t2 = Á 1 - 2 ˜
c Ë c ¯
By using Binomial expansion theorem, we obtain
2l Ê v2 ˆ
t2 = ÁË1 + 2 ˜¯ [Neglecting the higher order terms] (v)
c 2c
The distance travelled by the beam in time t2 is given by
x2 = ct2
Ê v2 ˆ
x2 = 2l Á1 + 2 ˜ (vi)
Ë 2c ¯
With the help of Eqs. (iii) and (vi), we get the path difference
Ê v2 ˆ Ê v2 ˆ
Dx = x1 - x2 = 2l Á1 + 2 ˜ - 2l Á1 + 2 ˜
Ë c ¯ Ë 2c ¯
È v2 v2 ˘
= 2l Í1 + 2 - 1 - 2 ˙
Î c 2c ˚
lv 2
\ Dx =
c2
Because of the introduction of this path difference of the two beams, the interference pattern would be shifted
as,
lv 2
n= [ path difference = nl for constructive interference]
lc2
where, l is the wavelength of the light used.

If the apparatus is rotated through 90°, the reflected and transmitted beams get interchanged and the path
lv 2
difference of will be produced in the opposite direction. This way the total path difference between the
c2
2lv 2
interfering beams becomes 2 and the interference pattern would be shifted as,
c
2lv 2
n= (vii)
lc2
In the above Eq. (vii) on putting l = 11 m, velocity of earth v = 3 ¥ 104 m/sec, c = 3 ¥ 108 m/sec and
l = 5.5 ¥ 10–7 m, the expected fringe shift comes out to be
2lv 2 2 ¥ 11 ¥ (3 ¥ 104 ) 2
n= =
lc 2
5.5 ¥ 10 -7 ¥ (3 ¥ 108 ) 2
or n = 0.4
This displacement of fringe width could have been observed easily since the apparatus used in this experiment
was capable of observing a fringe shift as small as 0.01. However, experimentally no significant fringe shift
could be observed. This experiment was repeated at different places on the earth, at different times of the
day and different seasons of the year. However, no fringe shift was observed in any case. This negative result
observed by the experiments suggests that the medium or space in which light propagates is not moving
relative to earth.
12.3.1 explanation of negative results of Michelson-Morley experiment

Three explanations were given in order to explain the negative results of Michelson-Morley experiment.
12.3.1.1 Constancy of Speed of Light

If we accept that the measured speed of light is the same for all directions in every inertial frame, then this
could be one of the basis to explain the negative results of Michelson-Morley experiment. If it is so, then
the speeds c + v and c – v will simply be equal to c. Hence, the times t1 and t2 will be equal to 2l/c. Thus
Dx = x2 – x1 will be zero. So according to the principle of constancy of speed of light, there should not be any
shift of fringes. Thus, the principle of constancy of speed of light successfully explains the negative result
of Michelson-Morley experiment. However, under this situation, there would be no experimental evidence
to indicate the existence of absolute ether. This is because all the inertial frames would be equivalent for the
propagation of light as it does not depend on the motion of the observer.
12.3.1.2 Ether Drag Hypothesis

In order to explain the negative results of Michelson-Morley experiment, we consider that the ether is
carried by the bodies moving through it or in simple words ether is dragged by the bodies moving through
it. As per this hypothesis, ether is also carried by the earth with its own velocity. Hence, there would be no
relative motion between the ether and the earth. Therefore, no shift of fringes is obtained. However, this
is contrary to stellar aberrations and Fizeau convection coefficient, which are the experimentally observed
phenomena.
12.3.1.3 Lorentz FitzGerald Contraction Hypothesis

As per this hypothesis, all bodies moving with velocity v are contracted in the direction of motion
by a factor 1 - v 2 /c 2 . So if L0 be the length of a body at rest with respect to ether and L be its length
when the body is in motion with velocity v with respect to ether, then L = L0 1 - v 2 /c 2 . With this
2L 2 L (1 - v 2 /c 2 ) 2 L (1 - v 2 /2c 2 )
t1 = (1 + v 2 /c 2 ) = 0 (1 + v 2 /c 2 ) = 0 (1 + v 2 /c 2 ) . If we neglect the higher
c c c
2L
powers of v/c in view of v < c, then the time t1 = 0 (1 + v 2 /2c 2 ) = t2 . It means the time taken or path
c
traversed by the reflected and transmitted beams is the same. So no shift of fringes is observed. However, this
concept could not be applied for explaining the negative results of Michelson-Morley experiment when the
two arms of the interferometer are not equal.
In view of the above explanations, it is clear that there is no acceptable experimental basis for the idea of ether
or absolute frame of reference. However, Einstein in 1905 gave a revolutionary idea that the motion through
ether is not a meaningful concept rather the motion relative to a frame of reference has physical significance.
Here the frame of reference could be the earth’s surface, the sun, the centre of galaxy, etc; though in every
case we must specify it. This idea finally developed the theory of relativity.
12.4 postulates of special theory of relativity LO3

We simply list the two postulates of the special theory of relativity, as
(1) All the fundamental laws of physics retain the same form in all the inertial frames of reference.
(2) The velocity of light in free space is constant and is independent of the relative motion of the source
and the observer.
12.5 LOrenTz TransFOrmaTIOn LO3

The invariance of speed of light in all inertial frames implies that Galilean transformation equations are not
suitable. Therefore, we have to introduce new transformation equations which are consistent with the new
concept of the invariance of speed of light in free space. These transformation equations were derived by
Lorentz and are known as Lorentz transformation equations.
Let us consider two inertial frames F and F¢, as shown in Y F Y¢ F¢ P(x, y, z, t)
Fig. 12.5. Again consider two observers O and O¢ situated at (x ¢, y ¢, z ¢, t ¢)
the origin in the frame F and F¢, respectively. Two coordinate
systems coincide initially at t = t¢ = 0. Let a pulse of light be
v
generated at time t = 0 from the origin which spreads out in
the space and at the same time the frame F¢ starts moving
with constant velocity v along +ve direction of x-axis
O O¢ X, X ¢
relative to the frame F. This pulse reaches at point P, whose
coordinates of position and time are (x, y, z, t) and (x¢, y¢, z¢, t¢)
measured by the observer O and O¢, respectively. Z Z¢
Figure 12.5
Therefore, the transformation equations of x and x¢ can be written as,`

x¢ = k (x – vt) (i)
where, k is a constant of proportionality and is independent of x and t.
The inverse relation can be written as,
x = k (x¢ + vt¢) (ii)
Keeping in mind that the time t and t¢ are not equal, we put the value of x¢ from Eq. (i) in Eq. (ii) in order to get
x = k[k ( x - vt ) + vt ¢ ]
x
or = kx - kvt + vt ¢
k
x kx
or t¢ = - + kt
kv v
kx Ê 1ˆ
\ t ¢ = kt - Á 1 - 2 ˜
v Ë k ¯
kx Ê 1ˆ
\ t ¢ = kt - Á 1 - 2 ˜¯ (iii)
v Ë k
Now according to second postulate of special theory of relativity speed of light c remains constant. So the
velocity of pulse of light which spreads out from the common origin observed by observers O and O¢ should
be the same. Therefore,
x = ct ¸
˝ (iv)
x ¢ = ct ¢ ˛
By putting the values of x and x¢ from Eq. (iv) in Eqs. (i) and (ii), we have,
ct ¢ = k ( x - vt ) = k (ct - vt )
or ct ¢ = kt (c - v) (v)
and ct = k (ct ¢ + vt ¢ )
or ct = kt ¢ (c + v) (vi)
By multiplying Eqs. (v) and (vi), we get
c 2 tt ¢ = k 2 tt ¢ (c 2 - v 2 )
c2
k2 =
(c 2 - v 2 )
1
or k=± (vii)
1 - v 2 /c 2
1 v2
or =1-
k2 c2
1 v2
or 1- = (viii)
k2 c2
Using Eqs. (i), (iii), (vii) and (viii), we have
x - vt
x¢ = (ix)
1 - v 2 /c 2
kx Ê v 2 ˆ
t ¢ = kt - Á ˜
v Ë c2 ¯
kxv Ê xv ˆ
= kt - = k Át - 2 ˜
c 2 Ë c ¯
Ê xv ˆ
ÁË t - 2 ˜¯
c
or t¢ = (x)
1 - v2 / c2
y¢ = y and z¢ = z (xi)
Hence, the transformation equations become
Ê xv ˆ
x - vt ÁË t - 2 ˜¯
c
x¢ = , y ¢ = y, z ¢ = z and t ¢ =
1 - v 2 /c 2 1 - v 2 /c 2
Imagine if the frame F is moving with velocity v along the –ve direction of x-axis relative to frame F¢, then
we get transformation equations of the form
Ê vx ¢ ˆ
x ¢ + vt ¢ ÁË t ¢ + 2 ˜¯
c
x= , y = y ¢, z = z ¢ and t =
1 - v 2 /c 2 1 - v 2 /c 2
These equations are known as inverse Lorentz transformation equations.

If the speed of moving frame is much smaller than the velocity of light c, (i.e., v << c) then the Lorentz
transformation equations reduce to Galilean transformation equations.
12.6 LenGTh cOnTracTIOn LO4

In classical mechanics, the length of an object is independent of the velocity of the moving observer relative
to the object. However, on the basis of the theory of relativity, the length of an object depends upon the
velocity of the observer with respect to the object.
Let us consider two inertial systems F and F¢. The Y F Y¢ F¢

system F¢ is moving with velocity v relative to the v
x ¢2
system F along x-axis, as shown in Fig. 12.6. Let a rod
x ¢1
AB be at rest in moving system F¢ relative to the observer
O¢ and L0 be the length of the rod in this frame measured O¢
A B
O
by observer O¢ at any instant. This length L0 measured X, X ¢
from the system in which the rod is at rest is called x1
x2
proper length. So L0 will be given by
Z Z¢
L0 = x2¢ − x1¢ (i) Figure 12.6
where x1¢ and x2¢ are the coordinates of two ends of the rod at any instant. At the same time, the length of this
rod (say L) measured by an observer O in the stationary frame F is given by
L = x2 − x1 (ii)
where, x1 and x2 are the abscissae of the ends of the rod in the frame F.
Since it would be appropriate (accurate in crude manner) to measure the length by the observer O¢ in the
frame F¢, we use Lorentz transformation equations for getting the length L0. For this, we have
x1 - vt
x1¢ = (iii)
1 - v 2 /c 2
x2 - vt
x2¢ = (iv)
1 - v 2 /c 2
Substracting Eq. (iii) from Eq. (iv) we have,
x2 - x1
x2¢ - x1¢ =
1 - v 2 /c 2
L
or L0 = (v)
1 - v 2 /c 2
v2
L = L0 1 - (vi)
c2
From Eq. (vi), we see that L < L0. Thus, the length of the rod is reduced in the ratio 1 - v 2 /c 2 : 1 as measured
by the observer moving with velocity v with respect to the rod.
12.6.1 physical insight
Length contraction is also called Lorentz contraction or Lorentz-FitzGerald contraction. The contraction
takes place only in the direction parallel to the direction in which the observed body travels. For example, in
the present case the length contraction takes place in the x-direction only. This effect is negligible at everyday
speeds for standard objects and can be ignored for all regular purposes. However, the effect becomes dominant
as the magnitude of the velocity approaches the speed of light. So the length contraction is the phenomenon
which is usually noticeable at a substantial fraction of the speed of light.
12.7 TIme DILaTIOn LO4

The word dilation means to lengthen. Consider two coordinate systems F and F¢ such that F¢ is moving with
velocity v along the x-axis relative to F. Imagine a gun placed at the fixed position (x¢, y¢, z¢) in the frame F¢.
Suppose it fires two shots at time intervals t1¢ and t2¢ measured by observer O¢ in the frame F¢.
The time interval (t2¢ − t1¢) of two shots measured by the clock at rest in moving frame F¢ is called proper
time interval and is given by
t2¢ − t1¢ = t0 (i)
As the motion is relative, we may assume that F is moving with velocity –v along the +x-axis relative to F¢.
In the frame F, the observer O, which is at rest, observes these two shots at different times t1 and t2. The time
interval appears to him is given by
t2 − t1 = t (ii)
Since it would be appropriate (accurate in crude manner) to measure the time by the observer O in the frame
F, we use inverse Lorentz transformation equations for getting the time difference t. For this, we have
t1¢ + vx ¢ / c 2
t1 = (iii)
1 - v 2 /c 2
t2¢ + vx ¢ / c 2
t2 = (iv)
1 - v 2 /c 2
By using Eqs. (iii) and (iv), we get,
t = t2 − t1
t2¢ - t1¢
or t=
1 - v 2 /c 2
t0
t= (v)
1 - v 2 /c 2
1
Eq. (v) shows that t > t0, i.e., the time interval appears to be lengthened by a factor which is
1 - v 2 /c 2
observed by the observer O in frame F. This is known as time dilation.
12.7.1 physical insight
The time dilation is actually a difference of elapsed time between two events, when measured by the observers
that are moving relative to each other, i.e., they have relative motion; the same effect is observed when the
observers are definitely situated from a gravitational mass (or masses). The two clocks with the two observers
may be measured to tick at different rates, which arises neither from technical aspects of the clocks nor from
the propagation time of signals. This takes place due to the nature of space-time, which is such that the time
measured along different trajectories is affected by differences in either velocity or gravity, as each of these
affects the time in different ways.
12.8 addition of velocities LO5

The classical laws of addition of velocities need to be modified at very high velocities. Consider two frames
of references F and F¢ such that the frame F¢ is moving with a velocity v relative to F along x-axis. Suppose a
particle is moving relative to both the systems F and F¢. Let u and u¢ be the velocities of the particle measured
in frames F and F¢, respectively. Then the velocity components are given as
dx dy dz ¸
ux = , uy = and u z = , Ô
dt dt dt Ô (i)
˝
dx ¢ dy ¢ dz ¢ Ô
u x¢ = , u y¢ = and u z¢ = ,
dt ¢ dt ¢ dt ¢ ˛Ô
From inverse Lorentz transformations, we have
x ¢ + vt ¢ t ¢ + vx ¢/c 2
x= , y = y ¢, z = z ¢ and t = (ii)
1 - v 2 /c 2 1 - v 2 /c 2
By differentiating these equations, we get
vdx ¢
dt ¢ +
dx ¢ + vdt ¢ c2
dx = , dy = dy ¢, dz = dz ¢ and dt = (iii)
1 - v 2 /c 2 1 - v 2 /c 2
From Eqs. (i) and (iii), we have

dx ¢
+v u¢ + v
dx dx ¢ + vdt ¢
ux = = = dt ¢ = x
dt vdx ¢ v dx ¢ v
dt ¢ + 2 1+ 2 1 + 2 u x¢
c c dt ¢ c
u¢ + v
ux = x (iv)
vu ¢
1 + 2x
c
Similarly,
dy ¢
1 - v 2 /c 2
dy dy ¢ 1 - v /c
2 2
uy = = = dt ¢ [ dy = dy ¢ ]
dt vdx ¢ v dx ¢
dt ¢ + 2 1+ 2
c c dt ¢
u y¢ 1 - v 2 /c 2
uy = (v)
vu ¢
1 + 2x
c
Similarly,
u z¢ 1 - v 2 /c 2
uz = (vi)
vu ¢
1 + 2x
c
Eqs. (iv), (v) and (vi) represent the relativistic laws of addition of velocities whereas in classical mechanics
ux is simply represented by
ux = ux¢ + v
If ux¢ = c, i.e., if the light is emitted in the moving frame F¢ along its direction of motion relative to F, then
u x¢ + v c + v c (c + v )
ux = = = =c
vu ¢ vc (c + v )
1 + 2x 1 + 2
c c
Thus, from the above expression it is clear that the speed of light is the same in all inertial frames.
vu x¢
If ux¢ and v are smaller as compared to c, then can be neglected as compared to unity and ux becomes
c2
ux = ux¢ + v, the law of addition of velocity which is similar to the one in classical mechanics.
12.9 variation of Mass with velocity LO6

In classical mechanics, mass of a particle is considered to be a Y F Y´ F´ v
constant quantity and independent of its velocity. However, in
relativistic mechanics, like length and time, the mass also depends B1 B2
on its velocity. Consider two frames of references F and F¢ such
that F¢ is moving with a constant velocity v relative to F in the u –u
positive direction of X, as shown in Fig. 12.7. Suppose, two similar O O´ X, X´
elastic balls B1 and B2 each having the same mass m approach
each other in the frame F¢ with equal speeds (u and –u) and collide
Z Z´
with each other in such a way that they coalesce into one body. By
applying the law of conservation of linear momentum, we have Figure 12.7
Momentum of ball B1 + momentum of ball B2 = momentum of coalesced body

(mu) + (–mu) = 0
Thus, the coalesced body must be at rest in frame F¢. Now, let us consider the collision with regard to the
frame of reference F, where u1 and u2 are the velocities of the balls. Then, according to Lorentz velocity
transformations
u+v
u1 = (i)
1 + uv /c 2
-u + v
u2 = (ii)
1 - uv /c 2
After an inelastic collision, the coalesced body moves with the velocity of frame F¢, (as it remains at rest
in F¢). Thus, v is the observed velocity in frame F. Let mass of the ball B1 moving with velocity u1 is m1 and
that of ball B2 moving with velocity u2 is m2 in the frame of reference F. By applying conservation of linear
momentum, we have
m1u1 + m2u2 = (m1 + m2)v (iii)
On substituting u1 and u2 from Eqs. (i) and (ii) into Eq. (iii), we have
È u+v ˘ È -u + v ˘
m1 Í 2˙
+ m2 Í 2˙
= (m1 + m2 )v
Î1 + uv / c ˚ Î1 - uv / c ˚
È u+v ˘ È -u + v ˘
m1 Í 2˙
- m1v = m2 v - m2 Í 2˙
Î1 + uv / c ˚ Î1 - uv / c ˚
È u+v ˘ È -u + v ˘
m1 Í - v ˙ = m2 Ív -
Î1 + uv / c
2
˚ Î 1 - uv / c 2 ˙˚
È u (1 - v 2 /c 2 ) ˘ È u (1 - v 2 /c 2 ) ˘
m1 Í ˙ = m2Í 2 ˙
Î 1 + uv / c ˚
2
Î 1 - uv / c ˚
m1 1 + uv / c 2
or = (iv)
m2 1 - uv / c 2
Now from Eq. (i),

2
Ê u + vˆ
ÁË ˜
u12 c ¯ (1 - u 2 /c 2 )(1 - v 2 /c 2 )
1- =1- = (v)
c 2
Ê uv ˆ (1 + uv / c 2 ) 2
ÁË1 + 2 ˜¯
c
Similarly, we can write
u22 (1 - u 2 /c 2 )(1 - v 2 /c 2 )
1- = (vi)
c2 (1 - uv/c 2 ) 2
On dividing Eq. (vi) by Eq. (v) we have
1 - u22 /c 2 (1 + uv / c 2 ) 2
=
1 - u12 /c 2 (1 - uv / c 2 ) 2
1 - u22 /c 2 1 + uv / c 2 (vii)
=
1 - u12 /c 2 1 - uv / c 2
Thus, from Eqs.(iv) and (vii), we have,
m1 1 - u22 /c 2
=
m2 1 - u12 /c 2
m1È 1 - u12 /c 2 ˘ = m2 È 1 - u22 /c 2 ˘ (viii)

ÎÍ ˚˙ ÎÍ ˚˙
From Eq. (viii), it is clear that the left hand side and right hand side are independent of one another. This
result may be correct only if each is a constant. Therefore,
m1 È 1 - u12 /c 2 ˘ = m2 È 1 - u22 /c 2 ˘ = m0
ÎÍ ˚˙ ÎÍ ˚˙
where, m0 is the rest mass of the body.
m0
Thus, m1 = (ix)
1 - u12 /c 2
m0
and m2 = (x)
1 - u22 /c 2
In view of Eqs. (ix) and (x), we conclude that if m0 be the rest mass of the body then its mass m when it moves
at speed v will appear as
m0
m= (xi)
∞
1 - v /c
2 2
This is the relativistic formula for the variation of mass with velocity. m
If we substitute v = c in Eq. (xi), then m becomes •, which means an object travelling
with the velocity of light would acquire infinite mass. Thus, no material particle can have v/c 1
a velocity equal to or greater than the velocity of light. The variation of mass m with the
velocity v/c is graphically shown in Fig. 12.8. Figure 12.8
12.10 eInsTeIn’s Mass energy relation LO6

In classical mechanics, the mass of a moving particle is independent of its velocity. Contrary to this aspect,
by Einstein’s special theory of relativity, the mass of a moving body depends upon its velocity and is given by
m0
m= (i)
1 - v 2 /c 2
where v is the velocity of the body, m0 is its rest mass and c is the velocity of light. The increase in energy
of the particle by the applications of force may be defined in terms of work done which is the product of the
force and the displacement. According to Newton’s second law of motion, the rate of change of momentum
of the particle is equal to the force applied on it. Thus
d (mv)
F= (ii)
dt
If the particle is displaced a distance dx by the application of force F, the work done Fdx is stored as kinetic
energy (EK) in the body. Then
dW = dEK = Fdx (iii)
By using Eqs. (ii) and (iii), we get
d (mv)
dEK = dx
dt
dx
or dEK = d (mv) (iv)
dt
or dEK = v[vdm + mdv] (as dx/dt=v)
= v 2 dm + mvdv (v)
m0
But m=
1 - v 2 /c 2
m02 c 2
or m2 =
c2 - v2
or m 2 c 2 - m 2 v 2 = m02 c 2 (vi)
By differentiating Eq. (vi), we have,
2mdmc 2 - 2mdmv 2 - 2vdvm 2 = 0 (as m0 and c are constants)
or dmc 2 = v 2 dm + mvdv (vii)

By using Eqs. (v) and (vii), we get
dEK = dmc 2 (viii)

Thus, from Eq. (viii) it is clear that the change in kinetic energy is directly proportional to the change in
mass. If the body is at rest, its velocity will be zero and hence the change in kinetic energy will be zero.
Therefore, its mass will be m0. If the body moves with velocity v, then its mass becomes m and its kinetic
energy becomes EK. Therefore, by integrating equation (viii), we get
EK m
Ú dEK = c 2 Ú dm
0 m0
or EK = c 2 [m - m0 ] = mc 2 - m0 c 2
or E = mc 2 = EK + m0 c 2 (ix)
From Eq. (ix), it is clear that E = mc2 is the total energy. It is the sum of kinetic and rest mass energy.
That is,
E = mc2 (x)
This relation is called Einstein’s mass energy relation.
s UmmarY
The topics covered in this chapter are summarized below.

✦ Newton’s second law of motion did not give correct results when applied to the objects moving with
high speed comparable to the speed of light. Moreover, for two observers, which are in relative motion,
the same set of transformation equations cannot be used to transform the laws of mechanics and
electromagnetism from the frame of reference of one observer to the frame of reference of the other
observer. The efforts of finding solutions to such problems finally led to the development of special
theory of relativity by Einstein in 1905.
✦ The systems in which the law of inertia holds good are called inertial systems or inertial frames of
reference. However, a non-inertial system is the one in which Newton’s first law of motion does not
hold good. A frame of reference which is in accelerated motion with respect to an inertial frame is
called non-inertial frame of reference.
✦ If the coordinates of a point P are (x, y, z, t) in a frame of reference F and (x¢, y¢, z¢, t¢) in another frame
of reference F¢ which is moving with uniform velocity v relative to F, then x¢ = x – vxt, y¢= y – vyt, z¢=
z – vzt and t¢= t are known as Galilean transformation of position. If ux, uy, uz be the velocity components
of the particle in the frame of reference F and u¢x, u¢y, u¢z be the velocity components in the frame of
reference F¢, then u¢x= ux – vx, u¢y= uy – vy, uz = uz – vz represent Galilean transformation of velocity of
particle. Galilean transformation of acceleration of particle is a¢x = ax, a¢y = ay, a¢z = az.
✦ The existence of stationary ether (absolute frame of reference) can be confirmed if we can measure
the absolute velocity of earth with respect to this ether. For this, Michelson-Morley experiment was
conducted at different places on the earth, at different times of the day and in different seasons of the
year. However, no shift of fringes was obtained. Thus the motion of the earth through ether could not
be experimentally detected.
✦ Three separate explanations were given for the negative results of Michelson-Morley experiment.
These explanations were based on ether drag hypothesis, Lorentz Fitzgerald contraction hypothesis
and the constancy of speed of light hypothesis.
✦ Invariance of speed of light in all inertial frames implies that Galilean transformation equations are
not suitable. Therefore, another transformation, named Lorentz transformation was introduced. The
Lorentz transformation equations are consistent with the new concept of the invariance of speed of
x - vt t - vx / c 2
light in free space. The transformation equations x ¢ = , y ¢ = y, z ¢ = z , t ¢ =
are
1 - v 2 /c 2 1 - v2 / c2
called Lorentz transformation equations, where (x, y, z, t) are the coordinates in a frame of reference F
and (x¢, y¢, z¢, t¢) in another frame of reference F¢, which is moving with uniform velocity v relative to F.
✦ In classical mechanics, length of an object is independent of velocity of the moving observer relative to
the object. However, as per theory of relativity the lengths measured in a frame F, which is at rest, and in
frame F¢, which is moving with velocity v relative to F, are not the same. However, L = L0 1 - v 2 /c 2
, where L0 is the length when measured in the system F¢ in which the object is kept and L is the length
when measured in the frame F which is at rest.
✦ As per theory of relativity the time intervals measured in a frame F, which is at rest, and in frame
F¢, which is moving with velocity v relative to F, are not the same. However, t = t0 / 1 - v 2 /c 2 ,
where t0 is the time interval when measured in the system F¢ in which the clock is kept and t is the time
interval when measured in the frame F which is at rest.
✦ Classical laws of addition of velocities are modified at very high velocities. The addition of velocities
u x¢ + v u y¢ 1 - v 2 / c 2 u z¢ 1 - v 2 /c 2
in relativistic mechanics reads u x = , uy = , uz =
1 + vu x¢ /c 2
1 + vu x¢ /c 2 1 + vu x¢ /c 2
✦ In relativistic mechanics, like length and time, the mass also depends on its velocity. If the mass of the
body is m0 when it is at rest, then its mass m when it moves with velocity v becomes m = m0 / 1 - v 2 /c 2 .
✦ In classical physics, the force F acting on a body is defined as the time rate of change of momentum.
However, in relativistic mechanics, it is the time rate of change of relativistic momentum. In view of
this, the kinetic energy of a particle of mass m0 which acquires velocity v when a force F acts on it
through a distance x in time t, is given by EK = (g – 1)m0c2, where g = 1/ 1 - v 2 /c 2 . This equation
states that the increase in kinetic energy of a particle is due to an increase in its mass. The total energy
of a particle E = mc2 is equal to the sum of rest mass energy (m0c2) and the kinetic energy of the particle.
E = mc2 is called the Einstein mass energy relation. This relation simply represents the total energy of
the particle in relativistic mechanics.
✦ The separate laws for the conservation of mass and the conservation of energy are replaced in the
theory of relativity by a single law called the conservation of mass energy or the law of conservation
of total relativistic energy. According to this law, the total relativistic energy is invariant under Lorentz
transformation. It means for an isolated system the total relativistic energy is the same as the energy
observed from any inertial system.
✦ A simple and very useful relation between relativistic momentum p, rest mass energy m0c2 and the total
energy E is given as E2 = p2 c2 + m02 c4.
s olved e XamPles
E xamplE 1 Show that if the variation of mass with velocity is taken into account, the kinetic energy of a
particle of rest mass m0 and moving with velocity v is given by
ÈÊ -1/ 2 ˘
v2 ˆ
K = m0 c 2 ÍÁ1 - 2 ˜ - 1˙
ÍÎË c ¯ ˙˚
Solution We know,
m0
K = (m - m0 )c 2 and m =
1 - v2/ c2
Ê m0 ˆ
\ K= - m0 c 2
Á ˜
Ë 1 - v /c
2 2
¯
Ê 1 ˆ
= m0 - 1 c2
Á ˜
Ë 1 - v /c
2 2
¯
= m0c 2 [(1 - v 2 /c 2 ) 1/ 2
- 1]

E xamplE 2 Show that the relativistic form of Newton’s second law, when F is parallel to
 -3/ 2
  dv Ê v2 ˆ
v , is F = m0 Á1 - 2 ˜ .
dt Ë c ¯
Solution We know that
m0 dp d
m= , p = mv and F = = (mv)
1 - v /c
2 2 dt dt
d È m0 ˘ d È v ˘
F= Í v ˙ = m0 Í 2 2 1/ 2 ˙
dt Í 1 - v / c ˙
2 2 dt Î (1 - v /c ) ˚
Î ˚
È 1 dv ( -1/ 2)( -2v / c 2 ) dv ˘
= m0 Í +v ˙
Î (1 - v /c ) dt
2 2 1/ 2
(1 - v 2 /c 2 )3/ 2 dt ˚
dv È (1 - v 2 /c 2 ) (v 2 /c 2 ) ˘
= m0 Í + ˙
dt Î (1 - v 2 /c 2 )3/ 2 (1 - v 2 /c 2 )3/ 2 ˚
dv 1 ÈÊ v2 ˆ v2 ˘
= m0 2 2 3/ 2 Í Á
1- 2˜ + 2˙
dt (1 - v / c ) ÎË c ¯ c ˚
-3/ 2
dv Ê v2 ˆ
= m0ÁË1 - 2 ˜¯
dt c
 -3/ 2
 dv Ê v2 ˆ
\ F = m0 Á 1 - 2 ˜
dt Ë c ¯
E xamplE 3 Derive energy-momentum relation or prove the relation

E2 – c2p2 = m02c4.
Solution Momentum-energy relation: According to mass-energy relation, we have, energy E = mc2 and momentum p
= mv
E2 – c2p2 = (mc 2)2 – c2(mv)2
E2 – c2p2 = m2c4 – c 2 m2v2
But according to mass-velocity relation,
m0
m=
1 - v 2 /c 2
2 2
È m0 ˘ 4 È m0 ˘ 2 2
\ E 2 - c2 p2 = Í c -Í c v
2 2 ˙ 2 2˙
ÎÍ 1 - v /c ) ˚˙ ÎÍ 1 - v /c ˚˙
m02c 4 m02c 2v 2
= -
(1 - v /c ) (1 - v 2 /c 2 )
2 2
È c2 - v2 ˘
= m02c 2 Í 2 2˙
Î1 - v /c ˚
È 1 - v 2 /c 2 ˘
= m02c 4 Í 2 2˙
Î1 - v / c ˚
E 2 - c 2 p 2 = m02c 4
E = c 2 p 2 + m02c 4
which is the required relation for energy and momentum.
E xamplE 4 Show that x2 + y2 + z2 – c2t2 = x¢2 + y¢2 + z¢2 – c2t¢2 or x2 + y2 + z2 – c2t2 is invariant under Lorentz
transformation.
Solution According to Lorentz transformation, we have
x - vt
x¢ =
1 - v 2 /c 2
y¢ = y
z¢ = z
t - vx / c 2
and t¢ =
1 - v 2 /c 2
2
È x - vt ˘
2
Ê t - vx/c 2 ˆ
\ x¢ + y ¢ + z ¢ - c t ¢ = Í
2 2 2 2 2
+ y 2 + z 2 - c2 Á
˙ 2 2˜
ÍÎ 1 - v /c ˙˚ Ë 1 - v /c ¯
2 2
( x - vt ) 2 - c 2 (t - vx / c 2 ) 2
= + y2 + z2
1 - v 2 /c 2
x 2 + v 2t 2 - 2 xvt - c 2t 2 - v 2 x 2 /c 2 + 2 xvt
= + y2 + z2
1 - v 2 /c 2
x 2 (1 - v 2 /c 2 ) - c 2t 2 (1 - v 2 /c 2 )
= + y2 + z2
1 - v 2 /c 2
( x 2 - c 2t 2 )(1 - v 2 /c 2 )
= + y2 + z2
1 - v 2 /c 2
= x 2 - c 2t 2 + y 2 + z 2
= x 2 + y 2 + z 2 - c 2t 2
This shows that the quantity x2 + y2 + z2 – c2t2 is same in both frames of references. So, this is invariant under Lorentz
transformation.
E xamplE 5 An event occurs at x = 100 m, y = 10 m, z = 5 m and t = 1¥10–4 second in a reference frame F.

Calculate the coordinates of the event in a reference frame F¢ which is moving with velocity 2.7¥108 m/sec
with respect to the frame F along the common XX¢ axis using (a) Galilean transformation and (b) Lorentz
transformation.
Solution (a) According to Galilean transformation, we have
x¢ = x−vt, y¢= y, z¢= z and t¢= t
Substituting the given values, we get
x¢ = 100 - 2.7 ¥ 108 ¥ 1 ¥ 10 -4 = - 26900 m

y ¢ = y = 10 m
z¢ = z = 5 m
and t ¢ = t = 1 ¥ 10 -4 sec
Thus, the coordinates in frame F¢ are
x¢ = 26900 m, y¢ = 10 m, z¢ = 5 m and t = 10–4 sec.

(b) According to Lorentz transformation, we have
x - vt t - vx/c 2
x¢ = , t¢ =
1 - v 2 /c 2 1 - v 2 /c 2
Substituting the given values, we get
100 - 2.7 ¥ 108 ¥ 1 ¥ 10-4

x¢ = = - 61720 m
1 - {(2.7 ¥ 108 )/(3 ¥ 108 )}2
y ¢ = y = 10 m,
z¢ = z = 5 m
1 ¥ 10-4 - {(2.7 ¥ 108 ¥ 100)/(3 ¥ 108 )2 }
t¢ =
1 - {(2.7 ¥ 108 )/(3 ¥ 108 )}
t ¢ = 2.288 ¥ 10-4 sec
Thus, the coordinates in F¢ are
x¢ = –61720 m, y¢ = 10 m, z¢ = 5 m and t = 2.288 ¥ 10–4 sec
E xamplE 6 At what speed must a particle move for its mass to be four times its rest mass?
Solution Given m = 4mo. v = ?
Formula used is
m0 m0
m= or 4m0 =
1 - v /c
2 2
1 - v 2 /c 2
v2 1 v2 1 15
1- 2
= or =1- =
c 16 c2 16 16
or v 2.9 108 m/sec
E xamplE 7 With what velocity a particle should move so that its mass appears to increase by 20% of its rest
mass?
20
Solution Given m = m0 + m = 1.2 m0
100
Formula used is
m0 m0
m= or 1.2m0 =
1 - v /c 2 2
1 - v 2 /c 2
v2 1 v2 1
1- 2
= or 2
=1- = 0.3055
c 1.44 c 1.44
v = 0.553 c
E xamplE 8 Show that the momentum of a particle of rest mass m0 and kinetic energy KE is given by the
expression.
K E2
p= + 2m0 K E
c2

m0
m=
1 - v 2 /c 2
\ m 2 (1 - v 2 /c 2 ) = m02
v 2 m 2 - m02
or =
c2 m2
or m 2v 2c 2 = [(m - m0 )c 2 ][(m + m0 )c 2 ]
= [(m - m0 )c 2 ][(m - m0 )c 2 + 2m0c 2 ]
= K E [ K E + 2m0c 2 ]
= K E2 + 2m0c 2 K E
K E2
or m 2v 2 = + 2m0 K E
c2
K E2
\ p = mv = + 2 m0 K E
c2
E xamplE 9 Calculate the mass and speed of 2 MeV electron.

Solution The relativistic kinetic energy of the electron is
KE = mc 2 - m0c 2
m0c 2
= - m0c 2
1 - v 2 /c 2
È 1 ˘
= m0c 2 Í - 1˙
ÍÎ 1 - (v /c )
2 2
˙˚
If an electron acquires this energy by a potential V volts, then
È 1 ˘
m0c 2 Í - 1˙ = eV
ÎÍ 1 - (v /c )
2 2
˙˚
1 eV
=1+
1 - (v /c )
2 2 m0c 2
1.6 ¥ 10 -19 ¥ 2 ¥ 106

=1+
9.1 ¥ 10 -31 ¥ (3 ¥ 108 ) 2
3.2 ¥ 10 -13
=1+
8.19 ¥ 10 -14
= 1 + 3.90 = 4.90
v2 1
1- 2
= = 0.204
c 4.90
v2
1- = 0.0416
c2
This gives
v = 0.98 c
m0
Mass, m =
1 - v 2 /c 2
m0 m0
m= =
Ê 0.98 c ˆ
2 1 - 0.96
1- Á
Ë c ˜¯
m0
m= = 5m0 = 5 ¥ 9.1 ¥ 10 -31 kg
0.2
m = 45.5 ¥ 10 -31 kg
E xamplE 10 Depict that no signal can travel faster than the velocity of light.
Solution
u x¢ + v
ux =
vu ¢
1 + 2x
c
If u¢x = c and v = c, where c = speed of light
c+c 2c
\ ux = = =c
c◊ c 2
1+ 2
c
Thus, addition of any velocity to the velocity of light simply reproduces the velocity of light. Hence, it can be concluded
that no signal can travel faster than the velocity of light.
E xamplE 11 At what velocity will the mass of a body is 2.25 times its rest mass?
Solution The mass of a body moving with velocity v is given by
m0
m=
1 - v 2 /c 2
Given m = 2.25 m0
m0
2.25m0 =
1 - v 2 /c 2
v2 1 1
1- = =
c 2 (2.25) 2 5.0625
1 v2
1- = 2
5.0625 c
v2
= 0.8024
c2
or v = 2.68 ¥ 108 m/sec
E xamplE 12 If the kinetic energy of a body is double of its rest mass energy, calculate its velocity.
Solution KE = (m – m0) c2
Given that
KE = 2m0c 2
2m0c 2 = mc 2 - m0c 2
3m0c 2 = mc 2
m0
3m0 =
1 - v 2 /c 2
1 - v 2 /c 2 = 1/3
or 1 - v 2 /c 2 = 1/9 or 1 - (1/9) = v 2 /c 2
8
8/9 = v 2 /c 2 or v = c
9
2 2
\ v= c = 0.94 c
3
E xamplE 13 The mass of a moving electron is 11 times its rest mass. Calculate its kinetic energy and
momentum.
Solution m = 11m0
KE = E - E0 = mc 2 - m0c 2
= 11 m0c 2 - m0c 2
= 10 m0c 2
= 10 ¥ 9.1 ¥ 10-31 ¥ (3 ¥ 108 ) 2 J
= 9.1 ¥ 9 ¥ 10-14 J
= 8.2 ¥ 10-13 J
8.2 ¥ 10-13
= eV
1.6 ¥ 10-19
= 5.1 MeV
Also,
m0
m=
1 - v 2 /c 2
m0
11 m0 =
1 - v 2 /c 2
1 1
1 - v 2 /c 2 = or v 2 /c 2 = 1 -
121 121
v = 2.98 ¥ 108 m/sec.
\ p = mv
= 11 ¥ 9.1 ¥ 10-31 ¥ 2.98 ¥ 108
= 298.298 ¥ 10-23
= 2.98 ¥ 10 -21 N-sec
E xamplE 14 How fast must an electron move in order to have its mass equal the rest mass of the proton?
m0
Solution m=
1 - v 2 /c 2
Given: m = 1.67 ¥ 10–27 kg (rest mass of proton)

m0 = 9.1 ¥ 10–31 kg (rest mass of electron)
9.1 ¥ 10 -31 kg
1.67 ¥ 10 -27 kg =
1 - v 2 /c 2
1 - v 2 /c 2 = 5.45 ¥ 10 -4
v = 2.999 ¥ 108 m/s
E xamplE 15 Find the velocity of a 0.1 MeV electron according to classical and relativistic mechanics.
Solution Classical mechanics gives
KE = (1/ 2)mv 2
0.1 ¥ 106 ¥ 1.6 ¥ 10-19 = (1/ 2) ¥ 9.1 ¥ 10-31 v 2
0.0351 ¥ 1018 = v 2
v = 1.87 ¥ 108 m/sec
Relativistic mechanics gives

KE = (m - m0 )c 2
Ê m0 ˆ
KE = - m0 c 2
Á ˜
Ë 1 - v /c
2 2
¯
KE 1
= -1
m0c 2 1 - v 2 /c 2
KE 1
1+ 2
=
m0c 1 - v 2 /c 2
0.1 ¥ 106 ¥ 1.6 ¥ 10-19 1
1+ -31
=
9.1 ¥ 10 ¥ (3 ¥ 10 )8 2
1 - v2/ c2
1
1 + 0.195 =
1 - v 2 /c 2
1
1.195 =
1 - v 2 /c 2
1 - (v 2 /c 2 ) = 0.7
This gives v = 0.54c = 0.54 ¥ 3 ¥ 108 = 1.64 ¥ 108 m/sec
m0
E xamplE 16 Prove that (1/2) mv2, where m = does not equal to the kinetic energy of a particle
1 - v 2 /c 2
moving at relativistic velocity.
Solution Relativistic kinetic energy is given by
ÔÏ m0 ¸
- m0 ˝Ô c 2
KE = (m - m0 )c 2 = Ì
ÓÔ 1 - v /c
2 2
˛Ô
ÔÏ 1 ¸
- 1Ô˝ m0c 2
=Ì
ÓÔ 1 - v 2 2
c Ô˛
-1/ 2
If v << c, then Ê1 - v ˆ
2
v2 3 v4 v2
ÁË ˜ =1+ + +=1+ 2
c2 ¯ 2c 2
8c 4
2c
v2 v4
Since is very small, the term and higher order terms are neglected.
c2 c4
Ï v2 ¸
\ KE = Ì1 + - 1˝ m0c
2
2
Ó 2c ˛
KE = (1/ 2) m0v 2
As m π m0
Hence, (1/2) mv2 does not equal to the KE of a particle moving at relativistic velocity.
E xamplE 17 Kinetic energy of a particle is (i) 3 times (ii) equal to its rest mass energy. What is its velocity?
Solution Formula used is EK = (m – m0)c2
m0c2 = rest mass energy
(i) Given, Kinetic energy = 3 ¥ Rest mass energy
(m – m0)c2 = 3 ¥ m0c2
or m = 4m0
m0 m0
and m= or 4m0 =
1 - v /c
2 2
1 - v 2 /c 2
2
v2 Ê 1 ˆ v2 1 15
1- =Á ˜ or =1- =
c Ë 4¯ c2 16 16
15
or v= c
16
or v = 0.968c = 2.9 ¥ 108 m/sec
(ii) Kinetic energy = Rest mass energy
(m - m0 )c 2 = m0c 2
m = 2m0
m0 m0
and m= or 2m0 =
1 - v /c 2 2
1 - v 2 /c 2
v2 1 3 3
=1- = or v = c = 0.866c
c2 4 4 2
= 2.6 ¥ 108 m/sec
E xamplE 18 Show that the circle x2 + y2 = a2 in frame F appears to be an ellipse in frame F¢ which is moving
with velocity v relative to F.
Solution The equation of circle in a stationary frame is x2 + y2 = a2
x = x¢ 1 - v 2 /c 2 and y = y¢
Substituting these values in the equation of the circle,
Ê v2 ˆ
x¢ 2 Á1 - 2 ˜ + y ¢ 2 = a 2
Ë c ¯
x¢ 2 Ê v2 ˆ y¢2
2 Á
1 - 2 ˜ + 2 = 1.
a Ë c ¯ a
a2
Suppose b 2 =
1 - (v 2 /c 2 )
x¢2 y¢2
\ 1 (This is the equation of an ellipse)
b2 a2
E xamplE 19 Calculate the mass m and speed v of an electron having kinetic energy 1.5 MeV. [Take rest mass
of electron m0 = 9.11 ¥ 10–31 kg and velocity of light c = 3 ¥ 108 m/sec].
Solution The relativistic kinetic energy K = (m – m0)c2
K = 1.5 MeV = 1.5 ¥ 106 ¥ 1.6 ¥ 10–19 J
1.5 ¥ 106 ¥ 1.6 ¥ 10–19 = (m – 9.11 ¥ 10–31) (3 ¥ 108)2
1.5 ¥ 106 ¥ 1.6 ¥ 10 -19
m - 9.11 ¥ 10 -31 =
(3 ¥ 108 )2
m = 3.58 ¥ 10 -30 kg
m0
m=
1 - v 2 /c 2
or 1 - (v 2 /c 2 ) = (m0 /m) 2
2
Ê 9.11 ¥ 10 -31 ˆ
v = c 1 - (m0 /m) 2 = 3 ¥ 108 ¥ 1 - Á ˜
Ë 3.58 ¥ 10 -30 ¯
v = 2.9 ¥ 108 m/sec
E xamplE 20 What is the length of a metre stick moving parallel to its length when its mass is (3/2) times of
its rest mass?
m0
L = L0 1 - (v 2 /c 2 ) and m =
1 - (v 2 /c 2 )
\ L = L0 (m0 /m)
= L0 [1/ m /m0 ] = L0 [1/(3/ 2)] [ m / m0 = 3/ 2]
or L = (2/3) L0 = 0.67 L0
For L0 = 1 m, L = 0.67 ¥ 1 = 0.67 m
E xamplE 21 A circular lamina moves with its plane parallel to the x-y plane of a reference frame S at rest.
Assuming its motion to be along the axis of x (or y), calculate the velocity at which surface area would appear
to be reduced to half to an observer in frame S.
Solution For an observer in frame S at rest, circular lamina, when in motion along the axis of x (or y), will appear to be
an ellipse. If diameter of the circle is D0, its value (say Dx) during motion will be
or Dx = D0 1 - v 2 /c 2
Area of elliptical lamina
ÊD ˆÊD ˆ
( Ae ) = p (a )(b) = p Á 0 ˜ Á 0 1 - v 2 /c 2 ˜
Ë 2 ¯Ë 2 ¯
2
ÊD ˆ v2
= pÁ 0˜ 1-
Ë 2¯ c2
Area of circular lamina
2
ÊD ˆ p D02
( Ac ) = p Á 0 ˜ =
Ë 2¯ 4
Given, Ae = Ac/2
p D02 v 2 Ê p D02 ˆ Ê 1 ˆ
\ 1- 2 = Á
4 c Ë 4 ˜¯ ÁË 2 ˜¯
Ê v2 ˆ 1 v2 1 3
ÁË1 - 2 ˜¯ = or =1- =
c 4 c2 4 4
3
v= c = 2.6 ¥ 108 m/sec
2
E xamplE 22 At what speed should a clock be moved so that it may appear to lose 1 minute in each hour?
Solution The clock loses 1 minute in 1 hour, means clock must record 59 minutes for each 1 hour. So that,
Proper time t0 = 59 min, apparent time t = 60 min.
According to Lorentz transformation, time dilation is given by
t0
t= (i)
1 - v 2 /c 2
Substituting the values in Eq. (i), we have

59
60 =
1 - v 2 /c 2
(1 - v 2 /c 2 ) = (59/60) 2
v 2 / c 2 = 1 - (59/60) 2
v 2 = [1 - (59/60) 2 ](3 ¥ 108 ) 2
v = 5.45 ¥ 107 m/sec
E xamplE 23 The proper life of p+-mesons is 2.5 ¥ 10–8 s. If a beam of these mesons of velocity 0.8 c is
produced, compute the distance the beam can travel before the flux of the meson beam is reduced to 1/e2
times the initial flux.
Solution
t0 2.5 ¥ 10 -8 2.5 ¥ 10 -8
t= = =
1 - v 2 /c 2 1 - (0.8c / c) 2 0.6
= 4.16 ¥ 10 -8 s.
If N0 is the initial flux and N is the flux after time t¢, we have
N = N 0e - t ¢ / t , where t is mean life time
\ N = (1/ e 2 ) N 0
\ N 0 / e 2 = N 0e - t ¢ / t
et ¢/ t = e 2
t ¢ /t = 2
t ¢ = 2t = 2t
The distance travelled by the beam before the flux is reduced to 1/e2 times the initial flux = 2t ¥ 0.8c = 2 ¥4.16 ¥10–8
¥0.8 ¥3 ¥108 = 19.96 m
E xamplE 24 A space ship moving away from the earth with velocity 0.6 c fires a rocket whose velocity
relative to the spaceship is 0.7c (i) away from the earth (ii) towards the earth. What will be the velocity of the
rocket, as observed from the earth in two cases?
Solution Formula used for relativistic velocity
u¢ + v
u=
vu ¢
1+ 2
c
where, u¢ is the velocity of rocket relative to space ship and v is the velocity of space ship relative to the earth.
The velocity away from the earth is taken as +ve and towards the earth as –ve.
Given u¢ = 0.7c and v = 0.6c.
(i) Rocket fired away from the earth, then
0.7c + 0.6c 1.3c
u= = = 0.915c
1 + 0.7 ¥ 0.6 1.42
= 0.92c
(ii) Rocket fired towards the earth, then
-u + v ( -0.7 + 0.6)c
u= = = - 0.17c
uv (1 - 0.7 ¥ 0.6)
= 0.17c
E xamplE 25 A 1.0 m long rod is moving along its length with a velocity 0.6c. Calculate its length as it
appears to an observer on the earth.
Solution Given L0 = 1.0 m and v = 0.6c.
Formula used is,
v2
L = L0 1 -
c2
2
Ê 0.6c ˆ
L =1 1- Á = 1 - (0.6)2
Ë c ˜¯
or
= 0.64 = 0.8 m
ExamplE 26 A rod has a length of 2.0 m. Find its length when it is carried in a rocket with a speed of 2.7 ¥ 108 m/sec.
Solution Given v = 2.7 ¥ 108 m/sec and L0 = 2.0 m.
Formula used is
v2
L = L0 1 -
c2
2
Ê 2.7 ¥ 108 ˆ
= 2.0 1 - Á ˜
Ë 3 ¥ 108 ¯
2
Ê 2.7 ˆ
= 2.0 1 - Á
Ë 3 ˜¯
= 2.0 1 - 0.81
= 2.0 0.19
= 0.872 m
E xamplE 27 Calculate the percentage contraction of a rod moving with a velocity Y
0.8 times the velocity of light in a direction at 60° to its own length.
Solution Given v = 0.8c.
L0 Sin 60°
L0
L0
In Fig. 12.9, component of length along the direction of motion Lx = L0 cos 60° =
2
3
and perpendicular to the direction of motion Ly = L0 sin 60° = L0 . 60°
2 X
The relativistic contraction occurs only along the direction of motion i.e., L0 Cos 60°
Figure 12.9
v 2 L0
Lx¢ = Lx 1 - = 1 - (0.8) 2 = 0.3L0
c2 2
and the component Ly along y-direction remains unchanged.
3
Ly¢ = Ly = L0 = 0.87 L0
2
The length of the rod in moving frame, i.e.,
L¢ = Lx¢ 2 + Ly¢ 2 = (0.3 L0 ) 2 + (0.87 L0 ) 2

= 0.196 L0
Percentage Length contraction

L0 - 0.916 L0
= ¥ 100
L0
= 8.4%
E xamplE 28 The length of a rod is found to be half of its length when at rest. What is the speed of rod
relative to the observer?
Solution Given L = L0 / 2, v =?
Formula used is
v2
L = L0 1 -
c2
L0 v2
or = L0 1 - 2
2 c
v2 1 v2 1 3
or 1- 2
= or 2
=1- =
c 4 c 4 4
v 3 3
or = or v = c
c 2 2
or v = 0.866 c
E xamplE 29 Calculate the length and orientation of rod of length 5m in a frame of reference which is
moving with a velocity 0.6c in a direction making an angle of 30° with the rod.
Solution Given v = 0.6c and L0 = 5m.
Refer to Fig. 12.9. The component of the length of the rod along x-direction will be
3
Lx¢ = L0 cos 30∞= 5
2
= 4.33
The component Ly remains unchanged, i.e.,
L0
Ly¢ = Ly = L0 sin 30∞ = = 2.5
2
The length of the rod in a moving frame (L¢), i.e.,
L¢ = Lx¢ 2 + Ly¢ 2 = (4.33) 2 + (2.5) 2 = 4.99 m
The orientation of the rod is given by angle q

Ly¢ 2.5
tan q = = = 0.57
Lx¢ 3.464
= 30.00°
E xamplE 30 Half-life of a particle at rest is 17.8 nanosecond. What will be the half-life when its speed is
0.8c?
Solution Given T0 = 17.8 ¥ 10–9 second and v = 0.8c

T0 17.8 ¥ 10-9 17.8 ¥ 10-9
T= = =
v 2
Ê 0.8c ˆ
2 1 - 0.64
1- 1- Á
c2 Ë c ¯˜
17.8 ¥ 10-9
= = 29.67 n sec
0.6
E xamplE 31 A clock keeps correct time on earth. It is put on the space ship moving uniformly with a speed
8
of 1¥10 m/sec. How many hours does it appear to lose per day?
Solution Given the time observed by the observer moving with the clock as 24 hrs, i.e., T = 24 hrs and the time observed
by the observer on the earth = T0
Formula used is
T0 T0 T0 2 2
T= = = fi T0 = 24 ¥
1 - v /c 2 2
Ê 1 ¥ 10 ˆ 8 2 8 3
1- Á ˜ 9
Ë 3 ¥ 108 ¯
or T0 = 8 ¥ 2 2 = 16 2
= 22.63 sec
Time lost per day = 24 – 22.63 = 1 hr 22 min 12 sec
E xamplE 32 With what velocity should a rocket move so that every year spent on it corresponds to 4 years
on earth?
Solution Given T0 = 1 year (proper time on the rocket) and
T = 4 year (relativistic time, corresponding time on the earth)
Formula used is
T0 1
T= or v =
1 - v 2 /c 2 1 - v 2 /c 2
v2 1 v2 1 15
or 1- = or =1- =
c 2 16 c2 16 16
15
or v= ¥ c = 0.968c
16
v = 0.97c
E xamplE 33 Determine the time (as measured by a clock at rest on the rocket) taken by a rocket to reach
a distant star and return to earth with a constant velocity v equal to 0.9999 c , if the distance to the star is 4
light years.
Solution Given distance of the star from the earth = 4 light year, T0 = time observed by observer in the rocket.
The time taken by the rocket to go to star from the earth and back with speed 0.9999 c is
2 ¥ 4c 8
T= year = = 8.0004 years
0.9999c 0.9999
T0
Now T= or T0 = T 1 - v 2 /c 2
1 - v /c
2 2
= 8.0004 1 - 0.9999 year

\ T0 = 0.08 year
E xamplE 34 In the laboratory, the life-time of particle moving with speed 2.8 ¥ 108m/sec is found to be
–7
2 ¥10 sec. Calculate the proper life-time of the particle.
Solution Given T = 2 ¥ 10–7 sec and v = 2.8 ¥ 108 m/sec.
Formula used is
2
T0 Ê 2.8 ˆ
T= or T0 = T 1 - v 2 /c 2 = 2 ¥ 10 -7 1 - Á
1 - v 2 /c 2 Ë 3 ˜¯
= 7.18 ¥ 10 -8 sec
E xamplE 35 Two electron beams travel along the same straight line but in opposite directions with velocities
v = 0.9c relative to the laboratory frame. Find the relative velocity of electrons according to Newtonian
mechanics. What will be the velocity measured by an observer moving with one of the electron beams?
Solution According to Newtonian mechanics, the relative velocity between electron beams will be
u¢ = u – v = 0.9c – (–0.9c) = 1.8c
The velocity measured by an observer moving with one of the electron beam
u - v 0.9c - ( -0.9c) 1.8c
u¢ = = =
uv
1- 2 1 + 0.9 ¥ 0.9 1.81
c
= 0.994 c
From the above, it is clear that the relative velocity according to Newtonian mechanics is found to be greater than the
velocity of light that is not possible.
E xamplE 36 Two photons approach each other, what is their relative velocity?
Solution Let velocity of each photon be c.
Formula used is
u¢ + v c+c
u= = =c
u ¢v 2
1+ 2 1+ c
c c2
i.e., the relative velocity of photons approaching each other is equal to the velocity of light.
E xamplE 37 A proton has a total relativistic energy as 900 MeV. If the rest mass of the proton is 1.6 ¥
10–27 kg, find its speed and kinetic energy.
Solution Given
E = 900 MeV = 900 ¥ 106 ¥ 1.6 ¥ 10 -19 J
= 1.44 ¥ 10 -10 J
Formula used is, E = mc2

E 1.44 ¥ 10 -10
\ Relativistic mass = m = = = 1.6 ¥ 10 -27 kg
c2 (3 ¥ 108 ) 2
From the above it is clear that
Rest mass of proton = Relativistic mass of proton.
So that proton is at rest and hence its speed and kinetic energy will be zero.
E xamplE 38 Dynamite liberates 5.4 ¥ 106 J/kg when it explodes. What fraction of total energy content is this?
Solution Total energy in one kg content = m0 c2
= 1 ¥ (3 ¥ 108)2 = 9 ¥ 1016 J
Dynamite liberates energy (per kg) = 5.4 ¥ 106 J
5.4 ¥ 106
Fraction of total energy = = 0.6 ¥ 10 -10
9 ¥ 1016
E xamplE 39 Calculate the speed of the electron which has kinetic energy as 1.02 MeV. Given rest mass
energy of the electron = 0.51 MeV.
Solution E = Kinetic energy + rest mass energy
= EK + m0c2
Given,
EK = 1.02 MeV = 2 ¥ 0.51 = 2m0c 2
= 2 ¥ rest mass energy
m0c 2
E = mc 2 =
1 - v 2 /c 2
m0c 2
\ EK + m0c 2 =
1 - v 2 /c 2
m0c 2
2m0c 2 + m0c 2 =
1 - v 2 /c 2
v2 1 v2 8
1- 2
= or =
c 9 c2 9
2 2
or v= c = 0.943c
3
= 2.83 ¥ 108 m/sec
E xamplE 40 The earth receives 1400 W/m2 of solar energy. The distance between the earth and the sun is
11
1.5 ¥ 10 m. Estimate the rate of decrease of the mass of the sun.
Solution Solar energy received by the earth = 1400 W/m2
= 1400 J/m2 sec
Distance of the earth from the sun R = 1.5 ¥ 1011 m
Total energy liberated by the sun per second
= 4p R 2 ¥ 1400 = 4 ¥ 3.14 ¥ (1.5 ¥ 1011 ) 2 ¥ 1400
= 3.96 ¥ 1026 J/sec
Energy at the sun is generated as per Einstein mass-energy relation, i.e.,

E = mc2 = 3.96 ¥ 1026 or m = 4.4 ¥ 109 kg/sec
E xamplE 41 Calculate the expected fringe-shift in a Michelson-Morley experiment if the distance from each
path is 2 meters and light has wavelength 6000 Å. Given, v = 3 ¥ 104 m/sec and c = 3 ¥ 108 m/sec.
Solution Given, l = 2 m, l = 6.0 ¥10–7 m, v = 3 ¥ 104 m/sec and c = 3 ¥ 108 m/sec.
The relation for fringe shift in Michelson-Morley Experiment is given by
2lv 2 2 ¥ 2 ¥ (3 ¥ 104 ) 2 4 ¥ 108
n= = -7
=
lc 2
6 ¥ 10 ¥ (3 ¥ 10 )8 2
6 ¥ 1016 ¥ 10
0.067
E xamplE 42 A clock is moving with a speed of 0.95c relative to an observer stationed on the earth. If the
speed is increased by 5% by what % does time dilation increases?
Solution Given, v = 0.95c and v ¢ = ÊÁ v +

5 ˆ
v˜ = 1.05 V
Ë 100 ¯
= 1.05 ¥ 0.95c
The relation of time dilation is given by
Dt Dt ¢ 1
Dt ¢ = or =
v2 Dt v2
1- 1-
c2 c2
Dt ¢ 1 1
= = = 3.2
Dt Ê 0.95c ˆ
2
1 - (0.95) 2
1- Á
Ë c ˜¯
Now v ¢ = 1.05 v = 1.05 = 0.95 c = 0.9975 c

1
Ê Dt ¢ ˆ 1
Then ÁË ˜ = = 14.14
Dt ¯ 1 - (0.9975)2
14.14 - 3.2
% time dilation increases = ¥ 100 = 77.3%
14.14
E xamplE 43 A beam of particle of half-life 2 ¥ 10–8 sec travels in the laboratory with speed 0.96 c. How
much distance does the beam travel before the number of particle is reduced to half-times of the initial value.
Solution The time interval in the particles own frame of reference in which the flux reduces to half of its initial flux is
the proper half-time (Dt ¢), given by formula
Dt 2 ¥ 10-8
Dt ¢ = = = 7.1 ¥ 10-8 sec
v2
1 - (0.96) 2
1-
c2
The distance travelled by the beam in this time in the laboratory frame
= 0.96c ¥ 7.1 ¥ 10–8 = 0.96 ¥ 3 ¥ 108 ¥ 7.1 ¥ 10–8
= 20.45 m
E xamplE 44 At what sped a body must move so as to have its mass double.
Solution Given, m = 2m0
Formula used is
m0
m=
v2
1- 2
c
m0
2m0 =
v2
1-
c2
v2 1 3
or =1- =
c2 4 4
3 3
v= c= ¥ 3 ¥ 108 m/ sec
2 2
= 2.6 108 m/sec
E xamplE 45 A muon decays with a mean life time of 22 ¥ 10–6 seconds measured in a frame of reference in
which it is at rest. If the muon velocity is 0.99c with respect to the laboratory, what is its mean life as observed
from laboratory frame?
Solution Given, Dt ¢ = 22 ¥ 10–6, v = 0.99c
Formula used is
Dt ¢ 22 ¥ 10-6
Dt = =
v2 1 - (0.99) 2
1-
c2
= 1.57 10 4
sec
E xamplE 46 A stationary body explodes into two fragments of rest mass 1 kg that move apart at speed of
0.6c relative to original body. Find the mass of original body.
Solution Given rest mass of two fragments (m0) = 1 kg and velocity of each fragment i.e., v1 = 0.6c and v2 = –0.6c
m0
Using the relation, m =
v2
1-
c2
m0 1 1
For fragment first, m1 = = =
v12 0.36c 2 0.64
1- 2
1-
c c2
= 1.25 kg
m0 1 1
For fragment second, m2 = = =
v22 0.36c 2 0.64
1- 1-
2
c c2
= 1.25 kg
By the law of conservation of mass, the mass of original body will be
M = m1 + m2 = (1.25 + 1.25) kg = 2.5 kg
E xamplE 47 What is the speed of particle whose KE is equal to its rest mass energy?
Solution Given, Rest mass energy = m0c2 and
Relavistic KE = (m – m0)c2
According to the problem, (m – m0)c2 = m0c2
mc2 = 2m0c2
m = 2m0
m0 v2 1
= 2m0 or 1- =
v2 c2 2
1-
c2
3
v= c or v = 2.6 ¥ 108 m/sec.
2
E xamplE 48 Find the energy equivalent to a mass of 5.0 mg.
Solution Given, Rest mass m0 = 5.0 mg = 5.0 ¥ 10–6 kg
As rest mass energy E = m0c2
E = m0c2 = 5 ¥ 10–6 ¥ (3 ¥ 108)2
= 45 ¥ 1010 Joules.
E xamplE 49 Show that E2/c2 – p2 is invariant.

Solution Consider two frames of reference F and F ¢. The frame F ¢ is moving with velocity v along x-axis with respect to F.
Let p be the momentum and E be the energy of a particle in reference frame F and p¢ and E ¢ in reference frame F ¢
\ p 2 = px2 + p 2y + pz2 (i)
p¢ =
2
px¢ 2 + p ¢y2 + pz¢ 2
1
and k= (ii)
Ê v2 ˆ
ÁË1 - 2 ˜¯
c
È V ˘
px¢ = k Í px - E2 ˙
Î c ˚
\ p ¢y = p y and pz¢ = pz
and E ¢ = k[E – vpx]
E¢2 2 k
2
- p ¢ = ( E - v px ) 2 - ( px¢ 2 + p ¢y2 + pz¢ 2 )
c2 c2
2
k2 È vE ˘
= 2
[ E - vpx ]2 - k 2 Í px - 2 ˙ - p 2y - pz2
c Î c ˚
k2 2 2 2 v2 E 2
= 2
[ E + v px¢ - 2vEpx - px2c 2 - 2 + 2vEpx ] - p 2y - pz2
c c
Èk Ï2
v E ¸
2 2 ˘
= Í ÌE 2 -
Î Ó
c 2
c 2 ˝
˛
{ }
+ - px2c 2 + v 2 px2 - p 2y - pz2 ˙
˚
Èk2 Ï Ê v2 ˆ Ê v2 ˆ ¸˘
= Í 2 Ì E 2 Á 1 - ˜ - px2c 2 Á1 - ˜ - p 2y - pz2 ˝˙
ÍÎ c ÔÓ Ë c ¯
2 Ë c ¯
2
Ô˛˙˚
v2 1
Put 1- =
c2 k 2
E¢2 k 2 È E 2 px 2c 2 ˘
2
- p¢ 2 = 2 Í - ˙ - p y - pz
2 2
c c Î k2 k2 ˚
E2
=
- ( px2 + p 2y + pz2 )
c2
E¢2 E2
2
- p¢ 2 = 2 - p 2
c c
E2
Hence, p 2 is invariant.
c2
E xamplE 50 A relativistic electron (m0 = 0.511 MeV/c2 and a photon (m0 = 0) both have momenta of 2.0
MeV/c. Find the total energy of each.
Solution Rest mass of electron m0 = 0.511 MeV/c2 and p(m0)photon = 0, pelectron = pphoton = 2.0 MeV/c
The momentum and energy relation for electron
E2 = e2p2 + m02c 4
(2.0) 2 (0.511)2 4
= c2 ◊ + ◊c
c2 c4
= (2.0)2 + (0.511)2 = 4.2611
E = 2.0642 MeV
The total energy for photon
MeV
E = cp = c ◊ 2.0 = 2.0 MeV
c
E xamplE 51 Show from Lorentz transformation that two events simultaneous (t1 = t2) at different positions
(x1 π x2) is a reference frame S are not, in general simultaneous in another reference.
Solution Consider a frame F ¢ moving relative to a frame F with a velocity v along x-axis. Let two event occur
simultaneously (t1 = t2) at different positions x1 and x2 (x1 π x2) in frame F and the corresponding times of occurrence in
frame F ¢ and t1¢ and t2¢ . According to Lorentz transformations,
x1 - vt1 x2 - vt2
x1¢ = and x2¢ =
Êv 2ˆ Ê v2 ˆ
1- Á 2˜ 1- Á 2˜
Ëc ¯ Ëc ¯
x1v x2v
t1 - t2 -
c2 c2
t1¢ = and t2¢ =
Ê v2 ˆ Ê v2 ˆ
1- Á 2˜ 1- Á 2˜
Ëc ¯ Ëc ¯
( x2 - x1 ) - v(t2 - t1 )
\ x2¢ - x1¢ =
Ê v2 ˆ
1- Á 2˜
Ëc ¯
For t1 = t2
x2 - x1
x2¢ - x1¢ =
Ê v2 ˆ
1- Á 2˜
Ëc ¯
v
(t2 - t1 ) - ( x2 - x1 )
Also t2¢ - t1¢ = c2
Ê v2 ˆ
1- Á 2˜
Ëc ¯
v ( x2¢ - x1¢ )
For t1 = t2, t2¢ - t ¢ = -
c2 Ê v2 ˆ
1- Á 2˜
Ëc ¯
v
or t2¢ - t1¢ = - 2 ( x2¢ - x1¢ )
c
Since x1¢ π x2¢ hence t1¢ π t2¢ . This shows that the two events which are simultaneous (t1 = t2) at positions (x1 π x2) in frame
F are not simultaneous in reference frame F ¢. The negative sign in above relations shows that the events occur at x2¢ .
First and then at x1¢ in frame F ¢.
E xamplE 52 A rod 1.0 m long is moving along its length with velocity 0.6c. Calculate the length as it
appears to an observer on the surface of earth.
Solution Let a rod be at rest in moving frame F ¢ relative to observer o ¢ and L0 be the lengths of the rod in this frame
i.e., L0 = 1.0 m.
L Ê v2 ˆ
Lo = or L = Lo Á1 - ˜
Ê v2 ˆ Ë c2 ¯
1- Á 2˜
Ëc ¯
2
Ê 0.6c ˆ
L = 100 cm ¥ 1 - Á = 100 1 - 0.36
Ë c ˜¯
= 100 0.64 = 80 cm
E xamplE 53 In an inertial F1 a red light and a blue light are separated by a distance Dx = 2.45 km, with the
red light at the longer value of x. The blue light flashes and 5.36 ms later the red light flashes. Frame F ¢ is
moving in the direction of increasing x with speed of v = 0.855c. What is the distance between the two flashes
and the time between them as measured in F ¢?
Solution Given: Dx = x2 – x1 = 2.45 km, Dt = t2 – t1 = 5.33 ms, v = 0.855c
Let (x1, t1) and (x2, t2) represent the position and time of blue and red light in frame F and ( x1¢, t1¢) and ( x2¢ , t2¢ ) are the
corresponding values in frame F¢
x - vt1 x - vt2
\ x1¢ = 1 , x2¢ = 2
2
v v2
1- 2 1- 2
c c
t1 - vx1 / c 2 t2 - vx2 / c 2
t1¢ = , t2¢ =
2
v v2
1- 1-
c2 c2
Let x ¢ and t ¢ represent the distance between the two flashes and the time between them, respectively, as measured in F ¢.
( x2 - x1 ) - v(t2 - t1 ) 2.54 ¥ 103 - 0.855 ¥ 3 ¥ 108 ¥ 5.35 ¥ 10-6
\ x¢ = x2¢ - x1¢ = = = 2.08 km
Ê v2 ˆ 1 - (0.855) 2
ÁË1 - 2 ˜¯
c
(t2 - t1 ) - v / c 2 ( x2 - x1 )
and t ¢ = t2¢ - t1¢ =
Ê v2 ˆ
ÁË1 - 2 ˜¯
c
-6 0.855
5.35 ¥ 10 - ¥ 2.45 ¥ 103
(3 ¥ 108 )2
= = 3.15 s
1 - (0.855)2
This result shows that when measurements are made from F ¢, the red flash comes before the blue flash in F ¢.
c
E xamplE 54 A particle of rest mass m0 moves with speed . What are mass, momentum, total energy and
2
kinetic energy.
c
Solution Given: m0 = rest mass, v =
2
m0 m0
By the relation m = =
Ê v2 ˆ Ê c ˆ
2
ÁË1 - 2 ˜¯ ÁË ˜
c 2¯
1- 2
c
m0
or m= = 2 m0
Ê 1ˆ
ÁË1 - ˜¯
2
m = 2 m0
c
Momentum p = mv = 2 m0 ¥ = m0c
2
or p = m0c
Total energy E = mc2 = 2 m0c 2
E = 2 m0c 2
and kinetic energy (KE) = (m – m0)c2 = ( 2 – 1)m0c2
or KE = (1.414 – 1)m0c2 = 0.414 m0c2
KE = 0.414 m0c 2
E xamplE 55 How fast must have an electron move in order to have its mass equal the rest mass of the proton?
Solution Given rest mass of proton (m) = 1.67 ¥ 10–27 kg
m0 = 9.1 ¥ 10–31 kg.
By using the relation

m0
m=
Ê v2 ˆ
1- Á 2˜
Ëc ¯
9.1 ¥ 10-31
1.67 ¥ 10-27 =
Ê v2 ˆ
Á1 - (3 ¥ 108 ) ˜
Ë ¯
v2
1- = 5.45 ¥ 10-4
(3 ¥ 108 ) 2
v2
= 29.69 ¥ 10-8
(3 ¥ 108 ) 2
v 2.999 108 m/sec.
E xamplE 56 Calculate the velocity that one atomic mass unit will have if its kinetic energy is equal to twice
the rest mass energy.
Solution Given kinetic energy = 2m0c2
As E = kinetic energy + rest mass energy
mc2 = KE + m0c2 [ E = mc2]
mc2 = 2m0c2 + m0c2 = 3m0c2
m
or = 3 or m = 3m0
m0
As m0 m0
m= or 3m0 =
È Ê v2 ˆ ˘ Ê v2 ˆ
Í1 - Á 2 ˜ ˙ Á
Ë
1- 2˜
Î Ë c ¯˚ c ¯
Ê v2 ˆ 1
ÁË1 - 2 ˜¯ = 9
c
v2 1 8
=1- =
c2 9 9
Ê 8ˆ
v = Á ˜ ◊c
Ë 9¯
v = 0.94 ¥ 3 ¥ 108
v 2.82 108 m/sec
E xamplE 57 A nucleus of mass m emits a gamma ray photon of frequency v. Show that the decrease in
È Ê hv ˆ ˘
internal energy of nucleus is not hv, but hv Í1 + Á 2˜˙
.
Î Ë 2mc ¯ ˚
Solution Given; frequency of gamma rays photon = v
The momentum of photon of frequency v is
hv
p=
c
hv
The nucleus of mass m recoils back with a momentum after emitting out a g-ray photon. The energy used in recoil is
c
2
Ê hv ˆ
1 2 p ÁË ˜¯ 2
(hv) 2
E = mv = = c =
2 2m 2m 2mc
Hence, the total decrease in internal energy of nucleus is given by
( hv )2 hv
hv hv 1
2mc 2 2mc 2
Hence proved.
E xamplE 58 What is the total energy of a 2.5 MeV electron?

Solution KE = 2.5 MeV
The total energy of a electron = KE + rest mass energy
m0c 2 9.1 ¥ 10-31 ¥ (3 ¥ 108 ) 2
E = 2.5 MeV + -19
= 2.5 MeV + eV
1.6 ¥ 10 1.6 ¥ 10-19
= 2.5 MeV + 0.512 ¥ 106 eV = 2.5 MeV + 0.512 MeV
= 3.012 MeV
E xamplE 59 Having the same momentum, which will move faster an electron or a photon?
Solution Given: pe = pp
or meve = mpvp
me >> mp, then
m
v p = e ve
mp
m
As me >> mp, then e >> 1, so
mp
vp >> ve
Photon will travel faster than electron.
E xamplE 60 Find the amount of work to be done to increase the speed of an electron from 0.6c to 0.8c. Take
rest energy of electron = 0.5 MeV.
Solution Given, m0c2 = 0.5 ¥ 106 eV
K = kinetic energy, E = K + m0c2
È m0 ˘
2Í
- m0 ˙
2 c 2 2
K = mc – m0c = Í v ˙
1- 2
ÍÎ c ˙˚
È 1 ˘
- 1˙
K = m0c 2 ÍÍ v 2
˙
1- 2
ÎÍ c ˚˙
È 1 ˘
Í - 1˙ 2È 1 ˘
K1 = m0c 2 Í Ê 0.6c ˆ
2
˙ = m0c ÍÎ 0.8 - 1˙˚
ÍÎ 1 - Á ˜ ˙˚
Ë c ¯
= 0.25 ¥ m0c2 = 0.25 ¥ 0.5 ¥ 106 eV
= 1.25 ¥ 105 eV
Similarly
2È 1 ˘
K2 = m0c Í - 1˙ = 0.67 ¥ 0.5 ¥ 106
Î 0.6 ˚
= 3.35 ¥ 105 eV
The amount of work done to be done = K2 – K1
= 2.1 ¥ 105 eV
= 2.1 ¥ 105 ¥ 1.6 ¥ 10–19 J
= 3.36 ¥ 10–14 Joule
3
E xamplE 61 What is the length of a meter stick moving parallel to its length when its mass is times of
2
its rest mass?
3
Solution Given: m = m0 , L0 = 1.0 meter
2
We know that
Ê v2 ˆ
L = L0 Á1 - ˜
Ë c2 ¯
m0
and m=
Ê v2 ˆ
ÁË1 - c ˜¯
2
m0 Ê v2 ˆ
fi = Á1 - ˜
m Ë c2 ¯
Ê m ˆ Ê 2ˆ È m 3˘
Or L = Lo Á 0 ˜ = Á ˜ Lo Í m = 2 ˙
Ë m ¯ Ë 3¯ Î 0 ˚
L = 0.67 ¥ 1.0 = 0.67 meter
L = 0.67 meter
E xamplE 62 How fast would a rocket ship have to go relative to an observer for its length to be contracted
to 99 per cent of its length at rest?
99
Solution Given l = lo
100
Ê v2 ˆ
Formula used l = lo Á1 - ˜
Ë c2 ¯
2
99 Ê v2 ˆ v 2 Ê 99 ˆ
lo = lo Á1 - ˜ or 1 – =Á ˜
100 Ë c2 ¯ c 2 Ë 100 ¯
v2 = 0.0199c = 0.0199 ¥ (3 ¥ 108)2
v= (0.0199) 3 ¥ 108
v 42.3 106 m/sec
E xamplE 63 A muon decays with a mean life time of 22 ¥ 10–6 seconds measured in a frame of reference in
which it is at rest. If the muon velocity is 0.99c with respect to the laboratory, what is its mean life as observed
from laboratory frame?
Solution Given: Proper mean life-time t0 = 22 ¥ 10–6 sec, v = 0.99c
\ Apparent mean life-time
t0 22 ¥ 10-6
t= =
v2
È (0.99c)2 ˘
1- Í1 - ˙
c2 Î c2 ˚
22 ¥ 10-6
= = 1.57 ¥ 10-5 sec
[1 - (0.99)2 ]
t 1.57 10 5 sec
E xamplE 64 At what speed should a rocket be move so that it may appear to lose 1 minute in each hour.
Solution Let the clock loses 1 minute in 1 hour, means clock must record 59 min for each 1 hour. So that
Proper time t0 = 59 min., Apparent time t = 60 min. According to Lorentz transformation,
t0
t=
Ê v2 ˆ
ÁË1 - 2 ˜¯
c
Substituting the values in above equation, we have
2
59 Ê 59 ˆ Ê v2 ˆ
60 = or Á ˜ = Á1 - ˜
v2 Ë 60 ¯ Ë c2 ¯
1-
c2
v2 Ê 59 ˆ
2
È 2˘
or =1- Á ˜ or v 2 = Í1 - ÁÊ 59 ˜ˆ ˙ ¥ 3 ¥ 108
c 2 Ë 60 ¯ Î Ë 60 ¯ ˚
v 5.45 107 m/sec
E xamplE 65 An electron has an initial speed of 1.4 ¥ 108 m/sec. How much additional energy must be
imported to it for its speed to double?
Solution Given: v = 1.4 ¥ 108 m/sec
Let rest mass of electron m0 = 9.1 ¥ 10–31 kg
The mass of electron at speed 1.4 ¥ 108 m/sec
m0
m=
v2
1- 2
c
9.1 ¥ 10-31
or m=
2
Ê 1.4 ¥ 108 ˆ
1- Á ˜
Ë 3 ¥ 108 ¯
= 1.029 ¥ 10-30 kg
Then corresponding energy

E1 = mc2 = 1.029 ¥ 10–30 ¥ (3 ¥ 108)2
= 9.261 ¥ 10–14 Joules
If v1 = 2v = 2 ¥ 1.4 ¥ 108 = 2.8 ¥ 108 m/sec
9.1 ¥ 10-31 9.1 ¥ 10-31
m¢ = =
Ê 2.8 ¥ 108 ˆ
2 0.359
1- Á ˜
Ë 3.0 ¥ 108 ¯
= 2.5348 ¥ 10-30 kg
E2 = mc2 = 2.5348 ¥ 10–30 ¥ (3 ¥ 108)2
= 22.813 ¥ 10–14 Joules
Hence, additional energy DE = E2 – E1
DE = 22.813 ¥ 10–14 – 9.26 ¥ 10–14
= 12.542 ¥ 10–14 Joules
DE = 1.254 ¥ 10–13 Joules
E xamplE 66 At what speed does a clock move if it runs at a rate which is one-third the rate of a clock at rest?
Solution Let a clock moves with velocity v w.r.t. another similar clock which is at rest. Let the time interval observed by
the clock at rest is t and the time interval observed by the clock moving with velocity v is t0.
The time interval measured by an observer that is stationary w.r.t. to the moving clock is proper time. From concept of
time interval
1
t = g t0 where g =
v2
1- 2
c
Given that t = 3 t0,
3t0 = g t0
g =3
Velocity of the moving clock is
1
=3
v2
1- 2
c
v = 0.9428c
E xamplE 67 At what speed does a meter stick move if its length is observed to shrink to 0.6 m?
Solution From the length contraction, when the observer and the object have relative motion, length appears to be
shortened as per the following relation
L0 1
L= where v =
g 1 - v 2 /c 2
Given that
L = 0.6 m; L0 = 1.0 m
This implies
g = 1.666
1
or = 1.666
v2
1- 2
c
v = 0.7998c
The speed of a meter stick is 0.7998c.
E xamplE 68 The average lifetime of a p meson in its own frame of reference is 26.0 ns. If the p meson
moves with speed 0.9c with respect to the Earth,
(a) What is its lifetime as measured by an observer at rest on Earth?
(b) What is the average distance it travels before decaying as measured by an observer at rest on Earth?
Solution
(a) The life time of p meson in its frame is 26.0 ns. This is the proper time observation as there is no relative motion
between observer and the object (say t0). The velocity of the p meson w.r.t. Earth is 0.9c.
Let the observed life time of p meson by an observer at rest on Earth be t. Then
1
t = g t0 where g =
v2
1- 2
c
26.0
t=
(0.9c) 2
1-
c2
26.0
t= ns
1 - 0.81
26.0
t= ns
0.19
t = 2.294 ¥ 26.0 ns
t = 59.64 ns
(b) Average distance it travels before decaying as measured by an observer at rest on Earth is given as L = lifetime
¥ velocity
L = 59.64 ¥ 10–9 ¥ 0.9 ¥ 3 ¥ 108
L = 16.10 m
E xamplE 69 Electrons in projection television sets are accelerated through a potential difference of 60 kV.
(a) Calculate the speed of the electrons using the relativistic form of kinetic energy assuming the
electrons start from rest.
(b) Calculate the speed of the electrons using the classical form of kinetic energy.
(c) Is the difference in speed significant in the design of this set?
Solution
(a) Total energy of a relativistic particles is given as E = KE + rest mass energy
KE = E – m0c2
KE = g m0c2 – m0c2 = (g – 1)m0c2
Rest mass energy of electron is m0c2

= 9.1 ¥ 10–31 ¥ 3 ¥ 108 ¥ 3 ¥ 108 = 81.9 ¥ 10–15 J
= 0.511 MeV = 511 keV
Kinetic energy, KE = (g – 1)m0c2
60 = (g – 1) ¥ 511 [Potential difference of 60 kV fi kinetic energy = 60 keV]
60 571
or g =1+ = = 1.1174
511 511
1
g = = 1.1174
v2
1- 2
c
or vrel = 1.3376 ¥ 108 = 0.446c
1 2
(b) From classical point of view kinetic energy = mv
2
1
60 ¥ 1000 ¥ 1.6 ¥ 10–19 = ¥ 9.1 ¥ 10–31 ¥ v2
2
60 ¥ 1000 ¥ 1.6 ¥ 10-19 ¥ 2
2
vclassical = = 21098.90 ¥ 1012 = 2.1098 ¥ 1016
9.1 ¥ 10-31
m
vclassical = 2.1098 ¥ 1016 = 1.4525 ¥ 108 = 0.484c
s
(c) The difference between the speeds of electrons accelerated through the potential of 60 kV is not significant in
case of TV as the distance travelled is of the order of 50 cm and the time lag thus produced is in of the order of
few nano-second. This can be neglected in present case.
E xamplE 70 Two powerless rockets are heading towards each other on a collision course. As measured by a
stationary observer at Earth, rocket A has speed 0.800c, rocket B has speed 0.600c, both rockets are 50.0 m
in length, and they are initially 2.52 Tm apart.
(a) What are their respective proper lengths?
(b) What is the length of each rocket as observed by a stationary observer in the other rocket?
(c) According to observer at earth, how long before the rockets collide?
(d) According to Rocket A, how long before they collide?
(e) According to Rocket B, how long before they collide?
(f) If the crew are able to evacuate their rockets safety within 50 min (their own time), will they be able
to do so before the collision?
Solution
(a) Given velocity of rocket A w.r.t. Earth = 0.8c
Length of rocket A w.r.t. Earth = 50 m
Velocity of rocket B w.r.t. Earth = 0.6 c
Length of rocket B w.r.t. Earth = 50 m
The distance between two rockets w.r.t. Earth = 2.52 Tm = 2.52 ¥ 1012 m.
The proper length of rocket A
L0
L=
g
1 1
L0 = g L = ¥L= ¥ 50 m
2
v (.8c) 2
1- 1-
c2 c2
1
L0 = ¥ 50 m = 83.33 m
0.6
Similarly, the proper length of rocket B
1 1
L0 = g L = ¥L= ¥ 50 m
2
v (.6c) 2
1- 1-
c2 c2
1
L0 = ¥ 50 m = 62.50 m
0.8
(b) To find the length of each rocket w.r.t. other rocket, we need to know the relative speed of each rocket w.r.t. each
other.
Thus, from relative speed formula, the speed of rocket A w.r.t. rocket B is
ux - v
u x¢ =
uv
1 - x2
c
0.8c - (-.6c)
or u x¢ =
.8c(-.6c)
1-
c2
1.4c
or u x¢ = ª 0.946c
1.48
Similarly the velocity of Rocket B as measured by a stationary observer in Rocket A is 0:946 c.
Then
L0
The length of rocket A w.r.t. rocket B is L =
g
83.33
or L= = 27.02 m
3.083
The length of rocket B w.r.t. rocket A is
62.5
L= = 20.27 m
3.083
(c) From Earth, the time taken in collision of two rockets is

0.8c ¥ Dt + 0.6c ¥ Dt = 2.52 ¥ 1012
2.52 ¥ 1012
Dt = = 6000 s = 100 min.
1.4 ¥ 3 ¥ 108
(d) The time taken before collision, according to rocket A

1
t= t0
v2
1- 2
c
1
100 = t0
(.8c) 2
1-
c2
(.8c) 2
t0 = 100 ¥ 1 - = 100 ¥ 0.6 = 60 min.
c2
(e) The time taken before collision, according to rocket B
1
t= t0
v2
1- 2
c
1
100 = t0
(.6c) 2
1-
c2
(.6c) 2
t0 = 100 ¥ 1 - = 100 ¥ 0.8 = 80 min.
c2
(f) All the crew members can evacuate their rockets before the collision as they get sufficient time for the same.
E xamplE 71 Superfast muons (v = .998c) can be produced by the collision of cosmic radiation with atoms
high in the atomosphere. Slow-moving muons in the laboratory fame have a lifetime as 2.2 ms. Experiments
show that a large number of muons do reach the sea surface. Explain this phenomenon with time dilation.
Solution The life time of muons is 2.2 ms in the laboratory frame i.e., this is the observation with no relative motion
between muons and observer. From time dilations, the life time of muons is different and increased from an observer
(Earth) with respect to which muons are moving.
Thus
1
t = g t0 where g =
v2
1-
c2
Putting t0 = 2.2 ms and v = 0.998c

1 1
t= t = 2.2 ms ({ d = ct)
2 0
v (.998c) 2
1- 2 1-
c c2
t = 15.81 ¥ 2.2 ms = 34.78 s
The distance covered by high speed muons during this life time calculated in Earth frame of reference is 10.41 km. This
results that a muon which is produced at height 10 km from sea level can travel the distance to reach at sea level with
speed .998c and life time 34.78 ms.
E xamplE 72 The period of a pendulum is measured to be 5 s in the reference frame of the pendulum. What
is the period when measured by an observer moving at a speed of 0.90c relative to the pendulum? What if we
increase the speed of the observer by 10%? Does the dilated time interval increase by 10% or more?
Solution The time period of the pendulum gets changed when the reference frame is changed from stationary frame to
moving frame. Thus
1
t = g t0 where g =
v2
1-
c2
Here t0 = 3s and v = 0.95c then
1
t= 3 = 2.294 ¥ 5s = 11.47 s
(.9c) 2
1-
c2
If the speed of observer is increased by 10%, then speed becomes
v = .9c ¥ 1.10 = 0.99c
With 10% increased velocity, the value of gamma factor is modified from 2.294 to 7.0888 which in turn, increases the
dilated time period of pendulum by over 200%.
E xamplE 73 A spacecraft is measured to be 150.0 m long and 30.0 m in diameter while at rest relative to an
observer. If this spacecraft now flies by the observer with a speed of 0.95c, what length and diameter does
the observer measure?
Solution From length contraction, the length observations made from moving frame are shortened by factor g . Here,
Actual length L0 = 150.0 m
Actual diameter d0 = 30.0 m
Let observations made by observer while spacecraft is moving with velocity 0.95c be L and d.
Then
L0 v2 (.95c) 2
L= = L0 1 - 2 = 120 1 - = 120 1 - .9025 = 120 0.0975
g c c2
L = 150 ¥ .3122 m
L = 46.84 m
The length of the space craft will appear to be 46.84 m if space
craft is moving with velocity 0.95c relative to observer.
However, diameter will be appear to be of 30 m as there is no
motion of spacecraft in the axis of diameter.
E xamplE 74 Imagine a motorcycle moving with a speed

0.70c past a stationary observer, as shown in adjoining
figure. If the rider tosses a ball in the forward direction with
a speed of 0.60c relative to himself, what is the speed of the
ball relative to the stationary observer?
Solution The speed of the motorcycle relative to the stationary
observer is v = 0.7c. The speed of the ball in the frame of reference
of the motorcyclist is u¢x = 0.6c. Therefore, the speed ux of the ball
relative to the stationary observer is
u x¢ + v
ux =
u¢ v
1 + x2
c
.6c + .7c 1.3c
ux = = = .9154c
.6c + .7c 1.42
1+
c2
The speed of the ball relative to the stationary observer is 0.9154c.
E xamplE 75 An electron, which has a mass of 9.11 ¥ 10–31 kg, moves with a speed of 0.850c. Find its
relativistic momentum and compare this value with the momentum calculated from the classical expression.
Solution Relativistic momentum of the electron is given as prel = mv = g m0v
Classical momentum of the electron is given as pclassical = m0v
1
prel = g m0v = m0v
v2
1- 2
c
1
prel = 9.11 ¥ 10-31 ¥ .85 ¥ 3 ¥ 108
(.85c) 2
1-
c2
prel = 44.0988 ¥ 10-23 = 4.049 10 22
kg.m/s
pclassical = 9.11 ¥ 10 ¥ .85 ¥ 3 ¥ 10 = 23.23 ¥ 10-23 = 2.323 10
-31 8 22
kg.m/s
Relativistic momentum is approximately 89% greater than the classical momentum.
E xamplE 76 An electron in a television picture tube typically moves with a speed u = 0.450c. Find its total
energy and kinetic energy in electron volts.
1
Solution Total energy of a relativistic particle E = g m0c 2 = m0c 2
v2
1- 2
c
1
E= ¥ 9.11 ¥ 10-31 ¥ 3 ¥ 3 ¥ 108 ¥ 108 J
(.45c) 2
1-
c2
E = 1.1197 ¥ 9.11 ¥ 10-31 ¥ 3 ¥ 3 ¥ 108 ¥ 108 J
1.1197 ¥ 9.11 ¥ 10-31 ¥ 3 ¥ 3 ¥ 108 ¥ 108
E= eV
1.6 ¥ 10-19
E ª 0.574 MeV
Kinetic energy of relativistic particle = total energy – rest mass energy
= (.574 – .511) MeV ({ rest mass energy = m0c2)
= 0.063 MeV
E xamplE 77 If the total energy of a proton is 2.5 times its rest energy, what is the speed of the proton?
Determine the kinetic energy of the proton in electron volts. What is the proton’s momentum?
Solution Total energy = 2.5 times rest mass energy
E = 2.5 m0c2;
E = g m0c2;
1
From above relations, g = 2.5 =
v2
1-
c2
v = .9165c = .9165 ¥ 3 ¥ 108 = 2.747 ¥ 108 m/s
Kinetic energy of proton = (g – 1) ¥ rest mass energy = (g – 1) m0c2
KE = (2.5 – 1) ¥ 1.67 ¥ 10–27 ¥ 3 ¥ 3 ¥ 108 ¥ 108 = 22.54 ¥ 10–11
= 2.254 ¥ 10–10 J
2.254 ¥ 10-10
KE = eV = 1.40875 ¥ 109 eV = 1408.75 MeV
1.6 ¥ 10-19
For relativistic momentum,
Using expression for total energy E2 = p2c2 + m02c4; E = 2.5m0c2
p2c2 = E2 – m02c4
p2c2 = 6.25m02c4 – m02c4 = 5.25m02c4
pc = 5.25 m0c 2 ; m0c 2 = 938 MeV
938
p = 5.25 MeV
c
MeV
p = 2149.2
c
E xamplE 78 A crew watches a movie that is 2 hours long in a spacecraft that is moving at high speed through
space. Will an Earthbound observer, who is watching the movie through a powerful telescope, measure the
duration of the movie to be (a) longer than, (b) shorter than, or (c) equal to 2 hours?
Solution The two events are the beginning and the end of the movie, both of which take place at rest with respect to
the space craft crew. Thus, the crew measures the proper time interval of 2 h. Any observer in motion with respect to the
spacecraft, which includes the observer on Earth, will measure a longer time interval due to time dilation.
E xamplE 79 Suppose astronauts are paid according to the amount of time they spend travelling in space.
After a long voyage travelling at a speed approaching c, would a crew rather be paid according to (a) an Earth-
based clock, (b) their spacecraft’s clock, or (c) either clock?
Solution (a) If their on-duty time is based on clocks that remain on the Earth, they will have larger paycheques. A shorter
time interval will have passed for the astronauts in their frame of reference than for their employer back on the Earth.
Q.1 The aim of Michelson-Morley experiment was to

(a) prove existence of ether (b) measure speed of light
(c) measure speed of earth relative to ether (d) test the isotropy of space
Q.2 The postulates of special theory of relativity are applicable to
(a) accelerated frames (b) inertial frames
(c) stationary frames (d) none of these
Q.3 Frame F is stationary and the frame F¢ is moving along the positive x-direction. A rod placed along
y-axis if observed from F frame appears
(a) contracted (b) elongated
(c) unchanged (d) nothing can be said
Q.4 A rocket is moving with a velocity 0.70 c. Velocity of the light with respect to the rocket is
(a) 0.7 c (b) c (c) 1.4 c (d) 0.35 c
Q.5 The relative velocity of two photons when they approach each other will be
(a) less than c (b) 0 (c) more than c (d) c
Q.6 The energy produced by one kg of mass, which is fully converted into energy, will be equal to
(a) 3 ¥ 1010 J (b) 9 ¥ 1016 J (c) 1018 J (d) 1 J
Q.7 A body of mass m falls through h meters. The decrease in its mass is equivalent to
(a) mgh/c2 (b) mgh (c) mghc2 (d) mgh/c
Q.8 At what velocity the kinetic energy of a body is equal to its rest mass energy
(a) 3c / 2 (b) c/2 (c) c/3 (d) 2c
Q.9 Relativistic transformations were suggested by
(a) Newton (b) Einstein (c) Huygens (d) Lorentz
Q.10 The apparent length of a meter rod moving parallel to its length with velocity 0.6c will be
(a) 0.8 m (b) 0.6 m (c) 1 m (d) 1.2 m
Q.11 When a body of rest mass 1 kg moves with velocity of light, its mass becomes
(a) 0 (b) • (c) 2 kg (d) 100 kg
Q.12 Einsten’s famous mass energy relation is
1
(a) E = m0c2 (b) E = mc2 (c) E = m c2 (d) none of these
2 0
Q.13 A rod of length L0 is kept in a frame F¢ which is moving with velocity of light in the direction of length.
The observed length of rod from a stationary frame of reference (earth) would be
(a) • (b) 0 (c) 10 L0 (d) 3 ¥ 108 L0
Q.14 The negative result of Michelson-Morley experiment was that
(a) it could not measure speed of light (b) it could not prove the existence of ether
(c) it could not show the shifting of fringes (d) it could not prove the electromagnetic nature of light
waves
Q.1 What do you understand by frame of reference?

Q.2 What are inertial frames of reference?
Q.3 Is Earth an inertial frame? If not why?
Q.4 What are Galilean transformations?
Q.5 Write Lorentz transformation equation.
Q.6 Write inverse Lorentz transformation equations.
Q.7 What are time dilation and length contraction?
Q.8 What are postulates of special theory of relativity?
Q.9 Explain the variation of mass with velocity?
Q.10 What do you understand by mass-energy equivalence relation?
Q.11 What is the rest mass of photon?
P ractice P roblems
general Questions
Q.1 Distinguish between inertial and non-inertial frames of references. Give one example of each. Is earth
an inertial frame? Give reasons.
Q.2 What is Newtonian principle of relativity? Discuss with examples. Why should laws of nature be the
same in all inertial frames of reference?
Q.3 What are Galilean transformations? Derive Galilean transformation equations for two inertial frames.
State and prove Galilean invariance.
Q.4 Prove that Newton’s law of motion are invariant under Galilean transformations?
Q.5 What are the quantities which are invariant under Galilean transformations?
Q.6 Show that a frame of reference having a uniform translatory motion (or moving with constant velocity)
relative to an inertial frame is also inertial.
Q.7 Show that the laws of conservation of momentum and energy are invariant to Galilean transformations.
Q.8 What was the objective of conducting the Michelson-Morley experiment? Describe the experiment.
How is the negative result of the experiment interpreted?
Q.9 What do you conclude from Michelson-Morley experiment? If ether does not exist in what medium
does light travel? What vibrates in light waves?
Q.10 (a) What efforts were made to explain the null results of Michelson-Morley experiment on the basis
of ether hypothesis?
(b) Draw the ray diagram in ether frame after 90° rotation of the apparatus.
Q.11 (a) Why the apparatus of Michelson-Morley experiment was rotated through 90°?
(b) Why did Michelson and Morley repeat the experiment during day and night and during all seasons
of the year?
Q.12 State and explain the fundamental (basic) postulates of special theory of relativity and derive Lorentz
space time trasformation equations on their basis.
Q.13 Derive Lorentz transformation equations for space and time coordinates and show that these equations
become the Galilean equations at very low speeds.
Q.14 Show by means of Lorentz transformation equations that
x¢2 – c2t¢2 = x2 – c2t2
Q.15 Derive Lorentz transformation equations and using them prove that moving clock appears to go slow.
Q.16 (a) On the basis of Lorentz transformations derive an expression for length contraction.
(b) Define proper length.
(c) A circle and a square are moving along x-axis. How will they appear to stationary observer?
Q.17 Apply Lorentz transformation to derive expression for length contraction and time dilation.
Q.18 What do you mean by length contraction at relativistic speed? Deduce the necessary expression for it?
Q.19 (a) What do you understand by time dilation? Establish a relation between proper and improper
interval of time.
(b) Give an example to show that time dilation is a real effect.
Q.20 What do you understand by time dilation? On the basis of Lorentz transformations discuss the variation
of time with velocity according to the special theory of relativity. Explain why does a moving clock
appear to run slow. Explain the terms, ‘proper time’ and ‘improper time’. Show that when v << c
Lorentz transformations for time reduce to Galilean transformations.
Q.21 Deduce an expression for variation of mass with velocity and depict it graphically. Also prove that no
material particle can have a velocity equal or greater than the velocity of light (c).
Q.22 Obtain the relativistic formula for the addition of velocities and also show that the velocity of light is
an absolute constant independent of the frame of reference.
Q.23 (a) Starting from Lorentz transformation equations for space and time co-ordinates derive equations
for relativistic addition of velocities. Hence, prove that no material particle can move with a
velocity greater than that of light.
(b) Show that the law agrees with velocity addition formula for non-relativistic velocities.
Q.24 Starting with Einstein’s velocity addition formula show that it is in conformity with principle of
constancy of speed of light.
m0
Q.25 (a) Derive the formula for relativistic variation of mass with velocity, i.e., m =
1 - v 2 /c 2
(b) Hence prove that it is not possible for a material particle to have a velocity equal to or greater than
the velocity of light.
Q.26 Obtain Einstein’s mass energy relation and discuss it. Give some evidence showing its validity.
Q.27 Establish mathematically Einstein’s mass energy relationship. Explain physical significance of this
relation. Mention nuclear phenomena supporting this relation.
Q.28 Write notes on the following
(i) Michelson-Morely experiment and its results (ii) Variation of mass with velocity
(iii) Lorentz-FitzGerald contraction (iv) Time dilation
(v) Mass-energy equivalence
Q.1 A space ship is 50 metre long on the ground, when it is in flight its length appears to be 49 metres to
an observer on the ground. Find the speed of the space ship. [Ans: 0.6 c]
Q.2 Calculate the percentage contraction of a rod moving with a velocity 0.8 times the velocity of light in
a direction inclined at 45° to its own length. [Ans: 17.5 %]
Q.3 The length of a rod is 100 m. If the length of this rod is measured by the observer moving parallel to
its length is 51 m, find the speed of the observer. [Ans: 0.86 c]
4 +
Q.4 A burst of 10 p mesons travels in a circular path of radius 20 m at a speed v = 0.99c. The proper mean
life of p+ meason is 2.5 ¥ 10–8 s.
(i) How many mesons would be left in a burst that had remained at rest at the origin for the same
period of time?
(ii) How many mesons survive when the burst returns to the point of origin?
[Ans: (i) N0p+ meson would survive, (ii) 920]
Q.5 Calculate the velocity at which electron mass is 3 times the rest mass. [Ans: 2.45 ◊ 108 m/sec]
Q.6 What should be the speed of electron so that its relativistic mass is twice its rest mass? [Ans: 0.87 c]
Q.7 Kinetic energy of a particle is twice its rest mass energy. What is its velocity? [Ans: 0.943 c]
Q.8 Calculate the velocity of 1.0 MeV electron. [Ans: 2.82 ◊ 108 m/sec]
Q.9 If one gram of a substance is fully converted into energy in one second, how many calories of heat will
be produced and how much power will be generated? [Ans: 9 ◊ 1013 J, 9 ◊ 107 MW]
Applied Nuclear Physics 13
Learning Objectives
decay, beta decay, gamma decay and
LO 1 Understand charge, mass, and size of nuclear radiation detectors
nucleus, angular momentum, magnetic,
LO 7 Know about discovery of neutron
electric, and statistical properties, parity
LO 8 Understand nuclear reactions along
LO 2 Know about charge independence and
with disintegration energy and
meson theory
threshold energy
LO 3 Learn about binding energy of nucleus
LO 9 Learn about nuclear fission, nuclear
LO 4 Discuss nuclear stability fusion and controlled fusion together
LO 5 Explain nuclear shell model and its with plasma, ignition temperature,
theory and applications and nuclear Lawson criterion and methods of fusion,
liquid drop model LO 10 Learn about particle accelerators
LO 6 Learn about radioactivity together with including linear accelerator, cyclotron,
laws of radioactive disintegration alpha betatron and plasma-based accelerators
Introduction
For many chemists, the atomic nucleus is nothing but a point charge, which carries most of the mass of the
atom. However, physicists’ perspective is different and it has been the field of research for them to investigate
how the protons and neutrons of the nucleus play important roles in the history and structure of the universe.
Actually, rapid progress in nuclear physics began after the discovery of the neutron in 1932. The discovery of
the neutron solved a known puzzle related to the spin of the nitrogen-14 nucleus, which was experimentally
measured as 1 basic unit of angular momentum, but at that time physicists could not find any way to arrange
21 particles (14 protons and 7 electrons of 14 7N) so as to give a spin of 1. However, the presence of the neutron
as an uncharged particle in the nucleus with spin ½ solved this problem. Moreover, the concept of neutron was
used to explain spin differences in many different nuclides (nuclide is an atomic nucleus as characterised by its
atomic number, its mass number and nuclear energy state). On the other hand, neutrons play an important
role in achieving energy through nuclear reactions. For example, nuclear fission is initiated by a slow neutron
and the resulting reactions produce 2.4 neutrons on an average. This way, a chain reaction is generated that al-
lows a self-sustaining mode of operation. However, neutrons are not effective initiators for light element fusion
reactors; the reason being their consumption during reaction rather than their production in fusion reactors.
Since no chain reaction takes place, applied fuel is needed to be added continually to keep the operation con-
tinuous. Now researchers have been investigating different means to achieve fusion, for example by using laser
field, accelerated particles, etc. Application of plasma has been a very attractive field and hence we discuss in
this chapter new topics such as plasma, Lawson criterion, fusion by inertial confinement, magnetic confinement,
laser fusion, etc. As mentioned particle acceleration also contributes to this field, various types of accelerators,
namely linear accelerator, cyclotron, betatron and plasma-based accelerators have been discussed. Other top-
ics on radioactive disintegration of nuclear, nucleus radiation and radiation detectors also have been covered.
13.1 Basic PrOPerties Of NucLeus LO1

According to Rutherford model, most of the mass of the atom is concentrated in a small spherical volume
called the nucleus, and the electrons are distributed around the nucleus. The radius of the nucleus is about
10–13 cm and the spherical volume where electrons are distributed carries radius of about 10–8 cm. The
protons and neutrons are the only constituents of the nucleus, which are called nucleons.
13.1.1 charge, Mass and size

A nucleus consists of Z number of protons each with charge +e, and N number of neutrons, which are neutral.
The total number of nucleons is called mass number A, i.e., A = Z + N. The charge of the nucleus is +Ze.
The rest mass of the proton, mp, is 1.00727663 amu, whereas the mass of the neutron, mn, is 1.0086654
amu. It is seen that the mass of the nucleus is not exactly equal to Zmp + Nmn, but it is less than this value.
The energy equivalent to this mass difference is known as binding energy, which is responsible for holding
neutrons and protons together in the nucleus.
The scattering of a-particles by the nucleus has been an important source of information with regard to the
nuclei, and in each case this showed departures from Coulomb scattering. Based on these experiments and
some theoretical models, it was deduced that the size of any atom is almost constant, but the size of any
nucleus depends on its mass number A. Clearly the size of the nucleus having large number of nucleons
should be large. Assuming a nucleus to be in the form of a sphere, its radius R is given by
R = r0A1/3
The value of r0 is found to depend on the type of experiment. In general, it is between 1.2 ¥ 10–13 cm and
1.48 ¥ 10–13 cm.
13.1.2 angular Momentum
Each nucleon in the nucleus is assumed to have both orbital and spinning motions just like an electron has
in an atom. It means each nucleon has both orbital and spin angular momenta. The magnitude of the spin
hÊ hˆ
angular momentum is Á h = ˜ . Its orientation in space can be described by only two states: the spin axis
2Ë 2p ¯
is either parallel or antiparallel to any given direction (say z-axis). So the component of spin along z-axis is
either h/2 or –h/2. In view of this, the total angular momentum i of each nucleon is
i = l ± s,
where l is the orbital angular momentum and s is the spin angular momentum. For nuclei having more than
one nucleon, l is replaced with L and s is replaced with S, which represent the corresponding total momentum
of all the nucleons. Hence, the total angular momentum of the nucleus is given by
I=L±S
I is actually a vector, whose magnitude is the maximum possible component in any given direction (z-axis), and
is represented by italic I. The value of I is an integral multiple of h for the nuclei with even mass numbers, and it
is an odd half-integral multiple of h for the nuclei with odd mass numbers. In particular, even-even nuclei (nuclei
with both Z and N even) carry zero value of I. This is also called that even-even nuclei have zero spin, where
nuclear spin refers to the total nuclear quantum number. The name nuclear spin which is frequently used for
total angular momentum of the nucleus is actually misleading. This incorrect usage was introduced before the
problem of the internal structure of nuclei had attained its present importance. Since then it has been continuing.
13.1.3 Magnetic Property
Magnetic property of a nucleus is associated with the nuclear magnetic moment. The motion of the nucleons
inside the nucleus should give rise to the nuclear magnetic moments like the electrons’ motion in the atom
provides the magnetic moment. This is true always unless the total angular momentum of the nucleus or
the nuclear spin is zero. If we assume spherically symmetric charge distribution of the nucleus, the nucleus
will give rise to a magnetic dipole moment only. The nuclear magnetic moments mI are measured in terms of
nuclear magnetons mN, given by mN = eh/2 mH. Here mH is the mass of hydrogen atom, which is equal to mass
eh m
of proton mp. Hence mN = = B . mB is known as Bohr magnetons, which is defined as the magnetic
2m p 1836
moment associated with an atomic electron in orbital motion with an angular momentum of 1h. This is given by
eh
mB = = 0.927 ¥ 10–23 J/Wb/m2. The value of nuclear magnetons mN is thus obtained as 5.05 ¥ 10–27 J/wb/m2.
2me
The measured values of mI are between –3mN and +10mN. When the magnetic moment of the nucleus is in the
opposite direction to the direction of nuclear spin, mI carries negative values. The positive valve of mI means
the directions of the magnetic moment of the nucleus is the same as that of the nuclear spin. The magnetic
moment of a proton is +2.79276 mN, whereas that of neutron is –1.191315 mN. This indicates that the proton
and neutron have a non-uniform charge distribution, which is also very complex.
The magnetic moment of a nucleus can also be represented in terms of nuclear gyromagnetic ratio gI and
nuclear g-factor gI, as mI = gIhI = gImNI.
Since I = 0 for nuclei containing even numbers of protons and neutrons, the even-even nuclei have no mag-
netic moment.
13.1.4 electric Property
Electric property of a nucleus is associated with the electric quadrupole moment, which is highly important
in connection with the shape of the nucleus. The electric quadrupole moment is a measure of the deviation
of the nucleus from its spherical symmetry. Under the situation of a deviation, the nucleus can be imagined
to be an ellipsoid of revolution with its diameter 2b along the axis of symmetry and diameter 2a along the
axis perpendicular to this. The quadrupole moment Q of the nucleus, when its electric charge is uniformly
distributed throughout the ellipsoid, is given by
2
Q = Z (b 2 - a 2 )
5
Clearly Q is zero for the nuclei having spherical symmetry (a = b) and whose charge is uniformly distributed.
As the formula suggests, the magnitude of electric quadrupole moment depends on the magnitude of nuclear
charge Z, size of the nucleus (magnitudes of b and a) and the extent of deviation (difference in b and a) from
spherical symmetry. The sign of Q may be positive or negative, depending on nature of Z and the values of
b and a. The Q value for a nucleus (deuterium) with one proton and one neutron is +0.00274 ¥ 10–24 cm2,
whereas an isotope of lutecium has a Q value of 7 ¥ 10–24 cm2 since it has 176 nucleons.
13.1.5 statistical Property
The concept of statistics is related to the behaviour of large number of particles. In general, the properties of
assemblies of electrons, protons, neutrons, photons and atomic nuclei cannot be described based on classical
statistics rather these follow the quantum statistics, i.e., Bose-Einstein statistics and Fermi-Dirac statistics.
In this connection, wave function is an important quantity that describes the particular system. A nucleon is
described by a function of its three space coordinates and the value of its spin. The wave function is said to
be anti-symmetric if it changes sign when three spatial and one spin coordinates of two identical particles
are interchanged, otherwise it is symmetric. Fermi-Dirac statistics apply to the system of particles which
are governed by anti-symmetric wave function. It also follows that the Pauli exclusion principle applies to
particles obeying the Fermi-Dirac statistics. The electrons, protons and neutrons obey this statistics as is
done by neuclei of odd mass number. On the other hand, all the neuclei with even mass number obey the
Bose-Einstein statistics. This statistics apply to the systems of particles which are governed by symmetric
wave function. Since nuclei with odd mass number have total angular momenta as odd half-integral multiples
of h and nuclei with even mass number have momenta as integral multiple of h, there is a direct correlation
between the total angular momentum of a nucleus and its statistics.
13.1.6 Parity
The wave function of a nucleus, to a good approximation, may be expressed as a product of two functions:
one of the space coordinates and the other depending only on the spin orientation. If the spatial part of its
wave function remains unchanged when the space coordinates (x, y, z) are replaced by (–x, –y, –z), the mo-
tion of the nucleus is said to have even parity. If the spatial part of the wave function changes sign on such
transformation of coordinates, the motion of the nucleus is said to have odd parity.
The parity of a nucleus in a given state is related to the orbital angular momentum L. If the value of L is even,
the parity is even and if the value of L is odd, the parity is odd.
13.2 NucLear fOrces LO2

The forces between nucleons, i.e., between proton
C Coulomb potential
and neutron, neutron and neutron, and neutron and 2
proton, are referred to as nuclear forces. These forces V(r) = K Ze (for proton
r
must be attractive; if it is not so, there would be no approaching
V(r ) nucleus)
stable nuclei. In view of Coulombic repulsion, these
D
forces should be much stronger than Coulomb forces R
O
between protons. Again if it is not so, there would B r
be no stable nuclei. Study of structures of atoms V(r) = 0 for neutron
and molecules reveals that the nuclear forces do not approaching nucleus
contribute to the formation or the properties of atoms
and molecules. In view of this, the nuclear forces must –V0 A
be short-range. These forces extend from the centre to (~ 40 MeV)
the surface of the nucleus or only a little beyond its Figure 13.1
surface. Figure 13.1 shows the cases when a proton
or a neutron is brought from an infinite distance toward a nucleus at O. The radius of the nucleus is R. When
the proton (thin line) is brought toward the nucleus, (case r > R), it experiences the potential V(r) varying as
1/r and the Coulomb force of repulsion varying as 1/r2. However, for this case (r > R), the neutron does not
feel any force and hence V(r) = 0 (thick line). For r < R, the proton as well as the neutron starts feeling the
attractive force of the nucleus. Since this attractive nuclear force is much stronger than Coulomb repulsive
force, it is represented by a negative potential –V0(M typically 40 MeV). The exact form of the potential inside
nucleus, i.e., between 0 and R is still not known.
13.2.1 charge independence
The charge independence of nuclear forces implies that the force between two protons (p-p), the force
between two neutrons (n-n) and the force between a neutron and a proton (n-p) are almost equal. The forces
are said to be charge symmetric, if only p-p M n-n. From the observation that N M Z M A/2 for the light and
medium weight stable nuclei for which Coulomb repulsion can be neglected, it is deduced that neutrons and
protons have a tendency to go in pairs. Examples of extraordinary stable nuclei are 42 He, 84 Be, 12 16
6 C, 8 O, etc.
Experimental evidences indicate that the nuclear forces show saturation.
13.2.2 Meson theory
In 1935, Japanese physicist Yukawa established a fact theoretically, which was suggested by Heisenberg
in 1932 that nuclear forces result from the constant exchange of massive particles between two nucleons.
Since this phenomenon is like the exchange of photons resulting in Coulomb force between two charged
particles, the massive particles taking part in nuclear forces were also called heavy quanta. According to
Yukawa, because of short range nature of nuclear forces, a nucleon is surrounded by a cloud of virtual
massive particles which are constantly emitted and absorbed by the nucleon; the same way as the electrical
charge is surrounded by a cloud of virtual photons. When a nucleon is brought near to another nucleon, a
particle emitted by one may be absorbed by the other or vice-versa. This way there is a constant transfer of
momentum from one nucleon to the other and hence a force is exerted between them.
The massive particle or the heavy quanta involved in nuclear exchange was given the name meson (p) which
can exist in three different forms. Its neutral form is called neutral pi meson or pion (p°), its negative form is
called negative pion (p–) and its positive form is called positive pion (p+). The exchange of neutral or charged
pions results in the interactions between the nucleons. In this process, neutron is converted into proton and
the proton is converted into neutron.
13.3 BiNdiNg eNergy Of NucLeus LO3
An atom denoted by ZA X is formed by bringing together Z protons, Z electrons and (A – Z) neutrons. In

principle, the mass of the atom should be equal to the sum of masses of Z protons, Z electrons and (A – Z)
neutrons. However, the mass M(A, Z) of the atom is less than this sum in the free state. This decrease in mass
DM is converted into energy DE = DMc2. This much energy is released in the process of formation of the
atom. In order to break the atom into its constituents, we will have to supply energy equivalent to this energy
released. This energy is called the binding energy BE.
The mass defect DM can be written as DM = Zmp + Zme + (A – Z) mn – M(A, Z)
Here mp, me and mn are the masses of the proton, electron and neutron, respectively. The mass mp + me can be
written as mH, i.e., mass of the hydrogen atom, if we neglect the small binding energy of the hydrogen atom.
Hence, the mass defect DM becomes

DM = ZmH + (A – Z) mn – M(A, Z)
= 1.0081437 Z + 1.0089830 (A – Z) – M(A, Z)
Since one amu corresponds to 931.145 MeV, the binding energy of the nucleus is written as
BE (in MeV) = 931.145 [1.0081437Z + 1.0089830 (A – Z) – M(A, Z)]
931.145
From this we obtain average BE per nucleon as BE/A = [1.0081437Z + 1.0089830 (A – Z) –
M(A, Z)] A
A graph of the binding energy per nucleon as a function of the mass numer is called binding energy curve,
which is shown in Fig. 13.2. With the exception of 42 He, 12 16
6 C and 8 O, the value of binding energy per nucleon
for almost all the nuclei lie on or close to the binding energy curve. From the curve, we notice some of the
outstanding features, as
9
7.6 MeV
8 8.8 MeV
6
BE/A (in MeV)
0 20 40 60 80 100 120 140 160 180 200 220 240

Mass number, A
Figure 13.2
(i) BE/A curve attains a flat maximum around A = 50 with the value ~8.8 MeV.
(ii) BE/A is low for nuclei with low A, but it increases very rapidly with increasing A (say upto A ~ 20).
(iii) Average BE/A in the region A ~ 20 and A ~ 160 is about 8.5 MeV and does not show much variation
(iv) BE/A decreases slowly and continuously for A>140 and it reaches a value of 7.6 MeV at A = 238 for
238
U.
13.3.1 Explanation
Consider a nucleus as the drop of liquid in which protons
and neutrons take the place of the molecules in the drop.
Two cases of nuclei are shown in Fig. 13.3 with small Light nucleus
number of nucleons (light nuclei) and sufficiently large Medium/heavy nucleus
number of nucleons (medium or heavy nuclei). This is Figure 13.3
clear from the figure that the nucleons, which are deep inside the nucleus, are attracted from all sides by
neighbouring nucleons. However, the nucleons setting on the surface are attracted from one side only or by
very less number of nucleons. This way the binding energy for the nucleons at the surface of the nucleus is
smaller than the binding energy of the nucleons inside the nucleus. For nuclei having smaller number of A,
there is a large fraction of nucleons at the surface. For this reason, BE/A is small for small A. On the other
side, BE/A for nuclei with large A is due to the Coulomb repulsive force between protons which reduces the
binding energy.
13.4 NucLear staBiLity LO4

In nuclei with smaller number of nucleons, Coulomb repulsion can be neglected in view of strong nuclear
forces. However, Coulomb repulsion becomes important for the nuclei having large number of nucleons.
Since this repulsion has the effect of making nuclei unstable, heavy nuclei tend to be made up of more number
of neutrons than protons. If we plot number of neutrons N versus number of protons Z for all known nuclei,
we notice a great deviation from the line N/Z = 1 in the case of heavy nuclei. The deviation reaches N/Z = 1.6
for A = 238. The curve N/Z = 1 is called the stability curve. The isotopes on both sides of the stability curve
are radioactive, which decay in such a way that the final product lies on the stability curve and is now stable.
The total binding energy of a nucleus depends not only on the ratio N/Z but also on whether these numbers
of neutrons and protons are odd or even. This is called odd-even effect. All the stable nuclei can be classified
into four groups, namely even-even, even-odd, odd-even and odd-odd, based on number of protons and
neutrons, respectively. Even-even nuclei having even number of mass number A are the most abundant, i.e.,
these nuclei are most stable, whereas odd-odd nuclei are very few. The stability of odd-even and even-odd
lies between the two extremes.
Data collected for stable nuclei suggests that nucleons tend to form neutron-proton pairs. This is called
pairing of nucleons, according to which nuclei that satisfy the condition A/2 = Z or A = 2Z are more strongly
bound together and any deviation from A = 2Z should decrease the binding energy.
13.5 NucLear MOdeLs LO5

In the absence of a detailed theory of nuclear structure, attempts were made to correlate nuclear data in
terms of various models. Several models were proposed. Each of them was based on a set of simplifying
assumptions and hence, was useful in a limited way only. The nuclear models include the shell or independent
particle model, liquid drop model, collective nuclear model and the optical model for nuclear reactions. The
shell model and the liquid drop models are the most important and useful models of nuclear structure.
13.6 NucLear sheLL MOdeL

If you look at the structure of the atom, quantum mechanics correctly predicted the closure or filling of
electronic shells. For example, the electronic shells for the elements He(Z = 2), Ne(Z = 10), Ar(Z = 18),
Kr(Z = 36) and Xe(Z = 54) are completely filled or closed. Moreover, the ionization potentials for these
elements are found to be higher than those of their neighbouring elements. The numbers 2, 10, 18, 36 and 54
are referred to as the atomic magic numbers of the periodic table. On the other hand, detailed investigations of
the curve of BE/A versus A and other properties of nuclei reveal that many nuclear properties vary periodically
in a sense similar to that of the periodic system of the elements. It was understood that the periodicities in the
properties of the nuclei might be due to a nuclear shell structure similar to the atomic shell structure. Most of
the nuclear properties show discontinuities near certain even values of Z or N. Experiments show that stable
nuclei result when either Z or N is equal to one of the numbers 2, 8, 20, 50, 82 and 126. These number are
called nuclear magic numbers, which have been interpreted as forming closed shells of neutrons or protons in
analogy with the filling of electron shells in atoms. The proton and neutron shells appear to be independent
of each other.
Following are some of the evidences for the nucleus having a shell structure and hence for the existence of
nuclear magic numbers.
(i) As the elements corresponding to atomic magic numbers were found to be chemically inactive, the
nuclei for which N or Z corresponded to nuclear magic numbers are found to be more stable than
their neighbours.
(ii) Nuclei with the above nuclear magic numbers have many more isotopes than their neighbours. For
example, Sn with Z = 50 has 10 stable isotopes whereas In with Z = 49 and Sb with Z = 51 each have
only 2 isotopes.
(iii) When BE/A obtained from nuclear disintegration data and mass spectrographic measurements is
plotted against A, the binding energy curve is found to have several kinks or breaks corresponding
to sudden increase in the value of BE/A. They have been found to occur for the nuclei corresponding
to the nuclear magic numbers.
(iv) Nuclei with magic numbers should have very low cross section, as the closed shells mean that there
is no more vacancy. The plot of neutron-absorption cross section versus number of neutrons (N)
indicated this to be true.
The above discussion indicates that neutrons and protons within the nucleus are arranged into shells within the
nucleus like electrons in atoms. Each shell is limited to a certain maximum number of neutrons or protons. The
resulting configuration is particularly stable and has an unusually low energy when a shell is filled or closed.
13.6.1.1 Theory
In the case of the structure of the atom we knew the form of Coulomb potential (proportional to 1/r) and
arrived at the system of electronic shells and magic numbers theoretically by solving the Schrödinger equation.
However, in the case of the structure of the nucleus we do not yet know the exact form of the nuclear potential
or the nuclear force, though investigations show that nuclear force are strongly attractive and they extend over
a very short range from the centre of the nucleus. In several theories it is assumed that each nucleon moves in
its orbit within the nucleus and its orbit is determined by a potential energy function V(r) which represent the
average effect of all interactions with other nucleons and is the same for each nucleon. Since each nucleon
is regarded as an independent particle which is governed by the potential V(r), nuclear shell model is also
called the independent particle model. The potential energy V(r) is analogous to the Coulomb energy but this
potential V(r) describing the nuclear attractions is quite different from the Coulomb potential. It has a form
kr 2
between the square well potential V = –V0 and the so called oscillator potential V(r) = –V0 + , where r is
2
the distance between the nucleon and the centre of force and k is a constant. The solution of the Schrödinger
equation with the square well potential and the oscillator potential does not give all the magic numbers.
Ê r2 ˆ
However, the combination of these two potentials, i.e., when V(r) = –V0 Á1 - 2 ˜ with R as the nucleus
Ë k ¯
radius, gives rise to all the magic numbers except magic number 28.
13.6.1.2 Applications
This model has been applied successfully to a variety of nuclear problems. For example, based on this it
could be possible to predict the total angular momenta of nuclei and that too in good agreement with the
experiments. So it is possible to assign values of the total angular momenta to nuclei for which the same has
not been measured, viz. b-radioactive nuclei. Based on the shell model a correlation between the distribution
of isomers and the magic numbers has also been found. Groupings or iselands of isomers are found just
below the magic numbers 50, 82 and 126, and there is a break at each of these numbers. It means isomerism
disappears when a shell is filled and does not appear again until the next shell is about half full. The shell model
when used to predict the total angular momenta of low-lying excited levels established that the conditions
for isomerism should exist below the magic numbers 50, 82 and 126 but not immediately above them. This
model also predicted correctly when isomerism should appear in the unfilled shells. Finally, the experimental
data on magnetic moments and electric quadrupole moments have also been interpreted in terms of this
model. For example, the quadrupole moment is zero or small at the proton numbers 2, 8, 20, 50 and 82.
13.7 NucLear Liquid drOP MOdeL LO5

The nature of variation of binding energy per nucleon curve was explained based on the concept of liquid
drop model of the nucleus. The nucleus was considered to be analogous to a drop of incompressible fluid
of very high density of about 1017 kg/m3. Saturation and short range properties of the nuclear force are
analogous to the properties of the forces which hold a liquid drop together. In this connection, mass and
binding energy of the nucleus become important for which a formula would be very useful. Such a formula
is called semi-empirical mass formula or semi-empirical binding energy formula. The main contribution to
this formula comes from various terms, namely volume energy, surface energy, symmetry effect, Coulomb
energy and odd-even effect of the nucleus.
Since the volume of the nucleus is proportional to the mass number A, the volume energy is given by EV = aV Å
where aV is a constant. The Coulomb energy between the protons tends to lower the binding energy. So its
effect should be a term with minus sign. This is given by
Z ( Z - 1)
EC = -4aC , where aC is a constant.
A1/3
This has already been seen that the binding energy is reduced because the nucleus has a surface. Since the
radius of nucleus will be proportional to A2/3, the surface energy is represented by ES = –aS A2/3, where as is
a constant.
For a given value of A, there is a particular value of Z for the most stable nuclide. The Coulomb effect is
small in the light nuclides, where the condition Z = A/2 or A = 2Z is applicable for the stable or symmetric
(spherical) nuclides. It means the departure from the symmetry or the term corresponding to the symmetry
effect would be proportional to some power of neutron excess (A – 2Z).
A detailed study of the symmetry effect shows that the term corresponding to this effect is given by
( A - 2Z )2
Er = - ar , where ar is a constant.
A
Finally odd-even effect can be represented by a term Ed whose value depends on Z and N. For nuclides
d d
containing even Z and even N, Ed = + whereas for the nuclides with odd Z and odd N, Ed = - .
2A 2A
By combinining all the terms we obtain the following formula for the binding energy BE
Z ( Z - 1) ( A - 2Z )2
BE = aV A – aS A2/3 – 4aC - ar + Ed
A1/3 A
The values of the constants aV, aS, aC and ar are determined by a combination of theoretical calculations and
adjustments to fit experimental values of the masses or binding energies. These are given by
aV = 14 MeV, aS = 13.1 MeV, aC = 0.146 MeV and ar = 19.4 MeV.
With these the semiempirical binding energy formula is obtained as
Z ( Z - 1)
BE (in MeV) = 14A – 13.1A2/3 – 0.584
A
( A - 2Z )2
–19.4 + Ed
A
Here Ed = 135/A for even A, even Z and Ed = –135/A for even A, odd Z.
For odd A, even Z and odd A, odd Z, Ed = 0.
13.8 radiOactivity LO6

Radioactivity is the disintegration of certain natural heavy elements such as radium, actinium, uranium and
thorium. The radioactive disintegration is a spontaneous process. This disintegration is accompanied by the
emission of a-rays (positively charged helium nuclei), b-rays (fast electrons) and g-rays (short X-rays). The
ultimate end product of the radioactive disintegration process is an isotope of lead. The radioactivity is of two
types, namely, artificial radioactivity and induced radioactivity.
13.8.1 artificial radioactivity
Artificial radioactivity is the radiation obtained from isotopes after high energy bombardment in an accelerator
(discussed later) by a-particles, protons and other light nuclei, or by neutrons in a nuclear reactor.
13.8.2 induced radioactivity
Induced radioactivity is the radioactivity induced in non-radioactive elements by neutrons in a reactor, or
protons or deuterons in a cyclotron or linear accelerator. X-rays or g-rays do not induce radioactivity unless
their energies are exceptionally high.
13.8.3 Laws of radioactive disintegration

In the process of radioactive disintegration, any nucleus disintegrates by emitting a particle or a g-ray or by
capturing an electron from the atomic shell. This way the original nucleus gets changed during the process
of radioactive decay. As mentioned, the radioactive disintegration is a spontaneous process. Depending upon
the type of the nucleus, the overall activity is a prolonged process varying from a few seconds to millions of
years. As disintegrations occur in a random manner, we can say that radioactive decay is a statistical process.
13.8.3.1 Disintegration or Decay Constant
Let us assume a constant probability l per unit time for the decay of each nucleus of a given element. Then
statistically, the probability of decay of each nucleus in time dt is ldt. If there are N undecayed nuclei at time
t, the number dN that will decay in time interval between t and t + dt is given by
dN = –ldtN (i)
where the negative sign shows that N decreases with increasing t. The probability constant l is called the
disintegration constant or decay constant. From Eq. (i) we have
dN
= - l dt
N
The above equation can be integrated under the limits of N and t. So we assume that initially there are N0
radioactive atoms at time t = 0. With this, we get
N t
dN
Ú N
= - Ú l dt
N0 0
Ê Nˆ
or ln Á = - lt
Ë N 0 ˜¯
or N = N 0 e - lt (ii)
where N is the number of atoms present at time t.
13.8.3.2 Activity
The activity A is nothing but the number of disintegrations per second of a sample. By knowing A, we can
detect the presence of a radioactive sample not by the radioactive atoms present but the radiation emitted by
these atoms when they disintegrate. So, the activity A is obtained from Eq. (ii) as
dN
A= = l N 0 e - lt = l N (iii)
dt
13.8.3.3 Half-Life Time
The half-life time, T1/2, of any sample is defined as the time interval in which the number of undecayed atoms
decreases by half. Clearly in the half-life time, the activity drops to half, i.e., to A0/2. If N0 be the number of
radioactive atoms at t = 0, then after t = T1/2, N0 becomes N0/2. So, from Eq. (ii) we obtain
N0
= N 0 e - lT1/ 2
2
ln 2 0.693
or T1/2 = = (iv)
l l
It is clear from this equation that the unit of disintegration constant or decay constant is 1/sec as the units of
T1/2 is sec.
13.8.3.4 Relation between Half-Life Time and Mean Life Time
Equation (ii) due to the exponential nature of decay states that a given radioactive sample takes an infinite
time to disintegrate completely. Individual radioactive atoms may have life times between zero and infinity.
Hence, it is meaningful to talk about the average or mean life time t, which is different from the half-life time
of a radionuclide. The mean life of a radioactive nuclide is defined as
t1dN1 + t2 dN 2 + t3 dN 3 + 
dt = (v)
dN1 + dN 2 + dN 3 + 
The above equation states that dN1 nuclei have life time t1, dN2 have life time t2, etc. If we consider dN to be
small, Eq. (v) can be written in the integral form by noting that dN1 + dN2 + dN3 + ....... = N0. So
N0 N0
Ú tdN Ú tdN
t= 0
N0
= 0
(vi)
N0
Ú dN
0
From Eq. (ii), we get
dN = - N 0 l e - lt dt
The limit of N from 0 to N0 corresponds to the limit of time t as • to 0 [from Eq. (ii)]. Hence
0
- Ú l tN 0 e - lt dt •
1
t= •
= Ú l te - lt dt =
N0 0
l
1
or t= (vii)
l
Therefore, the mean life time of an element is the reciprocal of its decay probability per unit time. Therefore,
using Eq. (iv), we get
T1/2
t= = 1.44 T1/2
0.693
This is the relation between mean life time and half-life time of a radioactive nuclide. Clearly the mean life
time is larger than the half-life time.
13.8.4 alpha decay
If a nucleus contains 210 or more nucleons, i.e., when nuclei are so large that the short-range nuclear forces
that hold them together are barely able to counterbalance the mutual repulsion of their protons, then the
process of alpha decay takes place in order to reduce the repulsion, or in other words, to increase the stability
of such nuclei by reducing their size.
Let us consider a parent nucleus X which disintegrates into daughter nucleus Y and an alpha particle in alpha
decay. Then
A
Z X Æ ZA-- 42 Y + 42 He (i)
4
In this process, 2He is called a particle. It is clear that with the emission of one alpha particle from the
nucleus, the atomic number is decreased by 2 and atomic weight is decreased by 4. Hence, because of
the different value of Z, the chemical nature of the daughter nucleus is different from the parent nucleus.
The condition for alpha decay can be obtained by applying the principle of conservation of energy and
linear momentum. Let Mp, md and ma be the rest masses of the parent nucleus, daughter nucleus and alpha
particles, respectively. Since initially the parent nucleus remains at rest before decay, its linear momentum
is zero. Therefore, the directions of daughter nucleus and alpha particle should be just opposite to conserve
momentum. Let Ti be the total energy before decay and Tf be the total energy after decay. According to the
law of conservation of energy,
or Ti = Tf
Mpc2 = mdc2 + Ud + mac2 + Ua, (ii)
where Ud and Ua are the nuclear energies of the daughter and alpha particle, given by
1 1
U d = md vd2 and U a = ma va2 .
2 2
Now Eq. (ii) can be written as
Ud + Ua = (Mp – md – ma)c2. (iii)
This equation represents total disintegration energy or Q-value and its value must be positive for spontaneous
emission. Hence, the condition for spontaneous alpha decay is that the rest mass of the parent nucleus must
be greater than the sum of the masses of the daughter and alpha particle.
In order to calculate the kinetic energy of the alpha particle, we use the conservation laws of linear momentum
and energy. Conservation of momentum yields
mdvd = mava (iv)
as the parent nucleus is at rest initially. Also,
1 1
Q = U d + Ua = md vd2 + ma va2
2 2 (v)
By substituting the value of vd from (iv) we get
2
1 Ê ma va ˆ 1 2
Q= md ÁË m ˜¯ + 2 ma va ,
2 d
Êm ˆ
or Q = U a Á a + 1˜ ,
Ë md ¯
Q
or Ua = .
1 + (ma / md )
From Eq. (i) ma/md @ 4/A – 4. Therefore, the kinetic energy of the a particle can be written as
A-4
Ua @ |Q|
A
From the above expression, it is clear that the alpha particle carries most of the disintegration energy, as A is
very large. The decay of nuclei by alpha emission cannot be explained classically but quantum mechanically
we can explain it. It is presumed that the parent nucleus before decay consists of the daughter nucleus and an
alpha particle. In the language of quantum mechanics, an alpha particle exists in one of the discrete energy
states of the daughter nucleus (say T0) to a spherical region and the potential barrier created by the daughter
nucleus restricts its motion. Classically, an alpha particle does not have enough energy to climb the barrier
but quantum mechanically the wave associated with the alpha particle has some probability to penetrate
the barrier. This effect is called quantum tunneling. The probability of finding the alpha particle one side of
the barrier is much less than to find it in the other side of the barrier, and inside the barrier the probability
decreases exponentially. This explains why alpha-particle emitters have long half lives for low T0.
13.8.5 Beta decay
Beta decay is a radioactive decay in which a beta particle that may be either an electron or a positron is
emitted. Since an electron or a positron cannot exist in the nucleus, we assume that this particle is created
at the time when the nucleus disintegrates. Electron emission takes place when a neutron is simultaneously
converted into a proton. Actually, three mechanisms are involved in beta decay. These are
1. b¯ decay (Electron or b¯ emission): In this mechanism, a nucleus decays by emission of an electron. It is
known as b¯ decay. Since b¯ is also called negatron, it is also sometimes referred to as negatron emission.
2. b+ decay (or Positron emission): In this mechanism, a nucleus decays by emission of a positron. It is
known as b+ decay.
3. Electron capture: In this mechanism, a nucleus decays by capturing an extra nuclear atomic electron.
The electron disappears because its mass is converted into energy.
The condition for spontaneous decay is obtained by using the principle of conservation of energy. By using
this, we can show whether a given unstable beta nucleus will decay by b– emission, b+ emission, or electron
capture. If we represent X and Y as the parent and daughter nuclei respectively, then the processes, which we
have discussed above, can be given as
A A 0
b– decay: Z XÆ Z +1Y + -1e (i)
b+ decay: A
Z XÆ A
Z -1Y + 0
+1e (ii)
A 0 A
Electron capture Z X + -1e Æ Z -1Y (iii)
The above equations are correct if we take the kinetic energy of the emitted beta particle as maximum. Let
Mp, md and me be the masses of X, Y and electron, respectively, and Ud and Uemax be the kinetic energies
of the daughter nucleus and the electron. Since the kinetic energy of X is zero, i.e., Up = 0 the principle of
conservation of energy yields
Mpc2 = mdc2 + Ud + mec2 + Uemax (iv)
The disintegration energy Q of this decay is given as
Q = Ud + Uemax = (Mp – md – me)c2 (v)
We can replace the nuclear masses Mp and md by the atomic masses M(Z) and M(Z + 1) by making use of relations
M(Z) = Mp + Zme
M(Z + 1) = md + (Z + 1)me (vi)
In view of these equations, Eq. (v) can be written as

Q = [M(Z) – M(Z + 1)]c2 (vii)
–
As for spontaneous decay, Q must be positive, the above equation states that b decay will occur only if the
atomic mass of the parent atom is greater than that of the daughter atom.
We can find the following Q-value for b+ decay by applying principle of conservation of energy to Eq. (ii)
Q = [M(Z) – M(Z – 1) – 2me]c2 (viii)
Similarly, for electron capture
Q = [M(Z) – M(Z – 1)]c2 (ix)
The mass number A of the nucleus does not change, i.e., DA = 0 in all the above mentioned processes, namely
b– emission, b+ emission and electron capture. So we can say that beta decay is an isobaric transformation.
In order to conserve charge in b– decay, a neutron is simultaneously converted into proton, when an electron
is emitted.
In the process of beta decay, if there are only two particles named daughter nucleus and beta particle (b– or b+),
this situation is similar to a two-body problem. The linear momentum is, however, not conserved in beta decay
because the observed directions of certain nuclides and electrons are almost never exactly opposite. Since the
daughter nucleus is very large compared to that of the beta particle, the beta particle should take away all the
disintegration energy. However, experimentally it is not observed, but is found that almost all the emitted beta
particles have energies less than the maximum energy or end-point energy, i.e., between 0 and Uemax. It means
the energy conservation does not take place in the process of beta decay. Moreover, the angular momentum is
also not conserved in beta decay because spin is not conserved in the nuclear reaction. This can be explained
with the help of spin of particles.
Let there be A nucleons before decay; the number remains same after decay also. Since the spin of each
nucleon is ½, in beta decay the spin is not conserved. For example,
n Æ p + e–
The spin of the neutron is ½, proton is ½ and also the spin of electron is ½. So from the above relation, we
observe that the spin is not conserved and hence angular momentum is also not conserved.
All the above discrepancies would be removed if an uncharged particle of small or zero rest mass and spin
½ is emitted in beta decay. In 1930, Pauli postulated the existence of a neutrino. The neutrino carries off
energy equal to the difference between Uemax and the actual kinetic energy. The momentum of a neutrino
also exactly balances those of the electron and the recoiling daughter nucleus. Actually in beta decay, two
types of neutrinos are involved. The first one is called neutrino (n) and the other one is called antineutrino
(n ) .
–
b–decay: n Æ p + b– + n (x)
b+decay: p Æ n + b+ + n (xi)
0
Electron capture: p+ –1 e Æn+n (xii)
Based on the above equations, we can conclude that to conserve energy, linear momentum and angular
momentum, the particle neutrino must have zero rest mass, zero charge but a spin of ½.
13.8.6 gamma decay
By emitting alpha, beta or other particles, the nucleus disintegrates and is usually left in the excited state. If
the excited nucleus does not have sufficient energy to emit another particle, like an excited atom, it returns
to its ground state by emitting photons. The energy of these photons is equal to the energy differences
between the various initial and final energy levels up to several MeV. These emitted photons from nuclei are
called gamma rays. A nucleus in the higher energy state Ti makes a transition to lower excited energy state
(or ground state) Tf , then the excess energy is governed by the equation
DT = Ti – Tf
This excess energy is emitted either by gamma-ray emission or by internal conversion. The process of
internal conversion is more frequent in a heavy excited nuclei. In this process, an excited nucleus is returned
into its ground state by giving up its excitation energy to one 27

12 Mg
of the orbital electrons, which is very close to it. The internal
conversion is an alternative of gamma ray emission. Both β−
these processes are caused by electromagnetic interaction. In
the gamma ray emission, the gamma ray spectrum consists
β−
of discrete energy levels. This indicates that the nucleus
has discrete energy levels. The relationship between decay
schemes and energy levels is shown in Fig. 13.4. This figure 1.015 MeV
shows the beta decay of 27 27
12Mg to 13Al. The half-life of the γ
decay is 9.5 min and it may take place in either of the two 0.834 MeV
27 27 *
excited states of 13Al. The resulting excited nucleus 13 Al then
undergoes one or two gamma decays to reach the ground state. γ γ γ
Let us suppose that an electron is in the K-shell. Then the

nucleus de-excites by giving its excitation energy to that 0
27
nearest K-shell electron. Such electrons are thereby knocked 13 AI
out of the atoms entirely. These electrons are called conversion Figure 13.4
electrons and the kinetic energy of the conversion electron is
given by
Ue = DT – UB
where DT is the energy with which gamma ray would have been emitted and UB is the atomic binding energy.
It means the emitted electron has an energy equal to the loss in nuclear kinetic energy minus the binding
energy of the electron. This process is similar to the photoelectric effect in which a nuclear photon is absorbed
by an atomic electron.
13.8.7 Nuclear-radiation detector
The detection of nuclear radiations depends upon their interaction with matter and especially on the excitation
and ionisation processes. Nuclear–radiation detector is a device in which the presence of radiation induces
physical change that is observable. In such studies, it is essential to know the exact number of incident
radiations and their energy at the detector. Different types of detectors have been used for detecting these
radiations. The most common detectors are discussed below in detail.
13.8.7.1 Ionisation Chamber
The ionisation chamber works on the principle that a
charged particle ionises the gas molecule when it is passed
through the gas. In the ionisation process, the number of
Rays
ion pairs are formed and they give the valuable information
about the nature of the radiation, i.e., whether it is an a or
b-particle, and energy of the incident particle.
The ionisation chamber consists of a cylindrical chamber
R To Electrometer
fitted with a pair of electrodes, which are kept at some
distance from each other, and a high potential difference is
applied between them (Figure 13.5.) The chamber is filled
Figure 13.5
with a gas like air or argon at normal pressure. One end of
the chamber is slated with a window, which is made of mica or nylon of thickness of about 0.002 mm and is
coated with graphite to make it conducting.
When a particle (or radiation) passes through the chamber, the gas in the chamber gets ionised and ions
are collected by opposite electrodes. They give rise to current to the external circuit, which is measured
by a meter. This ionisation current is proportional to number of pairs of ions and hence to the number of
radiations entering the chamber. The ionisation chamber is much less sensitive to b-particles in comparison
with a-particles. It is also too insensitive to g-rays because they do not produce enough ionisation. With the
help of this chamber we cannot count individual particles, but only the average effect of these particles can
be determined.
13.8.7.2 Geiger–Mueller Counter
Geiger–Mueller or GM counter is most efficient, accurate and useful device, which is used for detecting
individual particles such as a-, b-, g- and X-rays. It consists of a metallic cylindrical tube fitted with an
axial fine tungsten wire (Fig. 13.6). One end window of this tube is sealed by thin mica sheet through which
radiation can enter the tube. The whole arrangement is enclosed within a thin glass chamber. This tube is
filled with a gaseous mixture of about 90% argon and 10% ethyl alcohol at a pressure of 10 cm of Hg. The
potential of the order of 1000 volts is applied between anode and cathode. The value of this applied voltage
is adjusted to be somewhat below the breakdown potential of the gaseous mixture. When radiation enters the
GM tube, some of the argon atoms get ionised and produce number of ion-pairs and then electrons are moved
to the anode. Due to the shape of electrodes, the electrostatic field is radial and it acts strongly near the anode.
When electrons move towards the anode, they collide with gas molecules and produce further ionisation. In
this process, the multiplication of ions continues and as a result, avalanche of electrons is obtained. If the
exciting potential is sufficiently high, the secondary ionisation takes place and further another avalanche of
electrons is obtained. Thus, within no time almost entire volume of the gas in the tube is ionised and it leads
to an amplification as high as 108. In ionisation process, the total number of ions produced does not depend
upon radiation entered.
Cathode Tube
C
Anode Wire To scalar
Pulse Amplifier
OR
Ratemeter
Window R
Figure 13.6
Dead time The time during which the counter is incapable of responding to ions is known as dead time. The
reason of incapability of counter can be explained as follows. The electrons are collected rapidly on the anode
because of their light mass and leave behind a space charge of slow-moving positive ions. After the collection
of electrons, the space charge of positive ions becomes large enough and it is sufficient to cancel the applied
electric field and then further ionisation is stopped. The counter remains ionisation dead till space charge
of positive ions is collected by the cathode. After the removal of space charge, the counter again becomes
sensitive for further pulse recording. Actually, an electronic circuit is used to quench the discharge and pass
on an impulse to record the event.
13.8.7.3 Scintillation Counter
It is a very sensitive device used for detection and measurement of high energy nuclear radiations, viz, a-, b-,
and g-rays. These radiations are detected by means of fluorescence which they produce in certain materials
called scintillants. The selection of scintillant depends upon the radiation to be detected. For a-, b-, and
g-particles ZnS, naphthalene and NaI crystal, respectively, are used as scintillants.
The scintillation counter consists of scintillation chamber, photo-multiplier tube and electronic counter.
The scintillation chamber is made of an aluminium casing in which suitable scintillation crystal is placed.
When the radiation enters the crystal, it produces a tiny flash of light. This tiny flash of light is made to
fall on a transparent photosensitive layer of a photomultiplier tube where it ejects photoelectrons. These
photoelectrons are accelerated by the successive dynodes which are kept at progressively higher voltage.
The photoelectrons are pulled to the dynode 1 where a number of secondary electrons are emitted for each
primary photoelectron. These electrons are pulled to the dynode 2 where the electrons are further multiplied
by the secondary emission. Ultimately, the number of electrons is so much increased in successive stages that
a measurable current pulse is obtained. These pulses are sent to an electronic system where they are counted.
This counter has many advantages over a GM counter. The efficiency of this counter for counting g-rays is
comparatively much higher. The time of flight of the electrons through this tube is so small that it can count
about 106 particles per second.
13.8.7.4 Wilson’s Cloud Chamber: Visual Detector

Wilson in 1911 designed a cloud chamber, as shown in Fig. 13.7 to detect the radioactive particles and to
measure their energy. The dust-free air mixed with saturated vapour of water and any other liquid like alcohol
or ether is filled in a chamber fitted with a piston. a- or b-particles enter the chamber such that each particle
produces ions by ionising the air molecules on its path. Suddenly the piston is dropped down in order to reduce
the temperature by which the air is cooled and becomes supersaturated with water vapour. This way, the
vapour will condense in the form of drops on the ions along the path of particle. These tracks are illuminated
with light and photographed by camera. The track appears like a white line on a black background. Different
types of ionising particles produce different tracks.
Figure 13.7
The track of a-particles is thick, short and continuous. On the other hand,
the b-particles track is thin, dotted and long (Fig. 13.8). Thus, these α
particles can be distinguished on the basis of tracks.
13.8.7.5 Bubble Chamber: High Energy Detector
The cloud chamber is not suitable to detect highly energetic particles.
To overcome this problem, Glaser in 1952 invented the bubble chamber, β
which is almost inverse of cloud chamber. A bubble chamber mainly con-
sists of a heavy-walled pyrex bulb filled with a low-boiling-point liquid
like liquid-propane or liquid-hydrogen. This liquid is compressed by pass-
ing air through a pressure-regulating device and the proper temperature
around the pyrex bulb is maintained with the help of thermostat-controlled Figure 13.8
oil-bath.
It is well-known that the boiling point of the liquid can be raised by increasing pressure on its surface. If
the pressure on its surface is increased (like in a pressure cooker), the boiling does not start till a higher
temperature is reached. If the pressure is suddenly released then liquid becomes superheated. This superheated
state can be maintained for a few seconds. If an ionising particle passes through the liquid immediately after
the pressure is released, it leaves a trail of ions behind it. These ions, left in the track of a particle, act as
condensation centres and form vapour bubbles. This track of bubbles can be immediately illuminated and
photographed.
13.8.7.6 Semiconductor Detector
This device, shown in Fig. 13.9, is used as a particle
detector. It consists of a p-n junction which is
connected in reverse bias. The purpose of applying
reverse bias is to increase the thickness of the
depletion layer. This depletion region has no carriers
of either sign. When an ionising particle enters the
depletion region, the number of electron-hole pairs
are produced. Under the influence of applied reverse
bias, the electrons and holes are swept rapidly to the
+ve and –ve electrodes. Thus, it produces a current
pulse across the resistor R which is amplified and Figure 13.9
then counted.
13.9 discOvery Of NeutrON LO7

It was 1931 when Bothe and Becker in Germany observed that if very energetic a-particles (emitted from
Po) fell on certain light elements (Be, B, Li), an unusually penetrating radiation was produced. At first, it was
assumed to be g-radiation, but the experimental results were very difficult to be interpreted based on this. In
1932, Curie and Joliot in Paris showed that if this unknown radiation fell on paraffin or any other hydrogen-
containing compound, it ejects protons of very high energy. In the same year of 1932, based on a series of
experiments, the physicist James Chadwick showed that the new radiation consisted of uncharged particles
of approximately mass of the proton. These uncharged particles were called neutrons.
The discovery of the neutron was quite helpful in explaining a known puzzle involving the spin of the
nitrogen-14 nucleus (14N), which had been experimentally measured to be 1 basic unit of angular momentum.
Since 147N would be composed of 14 protons and 7 electrons and both protons and electrons carried an intrinsic
spin of half unit of angular momentum, there was no way to arrange these 21 particles to give a spin of 1.
Actually, all the possible pairings could give a net spin of half. However, after the discovery of neutron, 147N
was considered to consist of 3 pairs of protons and neutrons together with an additional unpaired neutron and
proton. Since an unpaired neutron and proton each contributed a spin of half in the same direction, the total
spin of 147N came out to be 1. After this, the concept of nuclear neutrons was used to explain spin differences
in many different nuclides. Finally, the neutron was accepted as a basic structural unit of atomic nuclei.
13.9.1 Neutron cross-section
The probability that a bombarding particle will interact in a certain way with a target particle is represented
in terms of nuclear reaction cross-section. Each target presents a certain area, called its cross-section, to the
incident particle. The incident particle interacts with the target if it is directed to this area.
Neutron absorption cross-section is the cross-section for a nuclear reaction which is initiated by neutrons.
For many materials, this rises to a large value at particular neutron energies due to resonance effects. For
example, a thin sheet of Cd forms an almost impenetrable barrier to thermal neutrons. The same way we can
define neutron scattering cross-section. Neutron cross-section is the measure of probability of scattering and
absorption of neutron when it approaches a nucleus. The incident neutron interacts with the nucleus if it is
directed to this area. The neutron cross-section is denoted by the symbol s and depends on the energy of the
incident neutron and size and nature of the target. This is expressed in barns (1 barn = 10–28 m2).
When a beam of mono-energetic neutrons is allowed to fall on a target slab having thickness dx, then its
intensity I is reduced to dI. This decrement in intensity depends upon incident intensity (I) and number of
nuclei per unit area of the slab. Let n be the number of nuclei per unit volume. Then ndx will be number of
nuclei per unit area. So that, we have,
dI µ I
µ ndx
\ dI µ Indx
or dI = –s Indx (i)
where proportionality constant s is known as the total cross-section of neutrons and the quantity sn is known
as attenuation coefficient. The negative sign is appearing as the intensity gets decreased with distance. Eq. (i)
can be written as
dI
= -s ndx, (ii)
I
Which on integration gives
I x
dI
Ú I
= - Ú s ndx
I0 0
or [ln I ]II0 = - s n[ x]0x fi ln I - ln I 0 = - s nx
Ê Iˆ
or ln e Á ˜ = - s nx (iii)
Ë I0 ¯
I
or = e -s nx (iv)
I0
I = I 0 e -s nx
This expression shows that the intensity of a transmitted beam gets decreased exponentially with increasing
thickness of the slab. On the other hand, the total cross-section s is calculated from Eq. (iii) as
1 Ê Iˆ
s=- ln
nx ÁË I 0 ˜¯
1 Ê I0 ˆ
s= ln Á ˜ (v)
nx Ë I ¯
13.10 NucLear reactiONs: cONservatiON Laws LO8

If a target material is bombarded by fast-moving particles such as protons, neutrons, electrons, deuterons or
alpha particles, then the target nuclei after the bombardment are usually different from what they were before.
The target nuclei may change their mass number or atomic number or both due to interaction with an incident
particle. This is called a transmutation and the reaction is called nuclear reaction.
The collision of a bombarding particle of rest mass mb and kinetic energy Ub with a target containing nuclei
of rest mass MT and kinetic energy UT results in the nuclear reaction.
Ub + NT Æ NR + y (i)
In this interaction a recoil nucleus NR of rest mass MR with kinetic energy UR and a light particle y of rest
mass my with kinetic energy Uy are emitted, as shown in Fig. 13.10. As per the law of conservation of energy,
total initial energy Ti that is the sum of the rest mass energies and kinetic energies must be equal to the total
final energy, Tf . That is
Ti = Tf
or mbc2 + Ub + MTc2 + UT = MRc2 + UR + myc2 + Uy (ii)
or [(UR + Uy) – (UT – Ub)] = [(MT + mb)c2 – (MR + my)c2] (iii)
or Uf – Ui = Mic2 – Mfc2 (iv)
Y Y vy
my
vb MT αy
X X
αR
mb
MR
VR
Before collision After collision

(a) (b)
Figure 13.10
where Ui = UT + Ub and Uf = UR + Uy together with Ui and Uf as the initial and final kinetic energies and Mi
and Mf as the total initial and final rest masses, given by Mi = MT + mb and Mf = MR + my. The above equation
states that the net increase in kinetic energy is equal to the net decrease in rest mass energy that is equal to the
disintegration energy or Q-value. Hence Q-value is defined as
Q = Uf – Ui = Mic2 – Mfc2
13.10.1 Exoergic and Endoergic Reactions

If the Q-value is positive, then the nuclear reaction is called exoergic reaction and in this case Uf > Ui or
Mic2 > Mfc2. If Q-value is a negative quantity, then the nuclear reaction is an endoergic reaction and in this
case Uf < Ui or Mic2 < Mfc2. If initially the target nucleus is at rest, then its kinetic energy UT = 0. Hence,
Q-value of a reaction is given by
Q = (UR + Uy) – Ub (v)

= [(MT + mb) – (MR + my)]c2 (vi)
13.10.2 disintegration energy
We can find the disintegration energy by knowing the values of Ub, Uy and UR. The kinetic energy UR of the
recoiling nucleus NR is very small. It is difficult to measure it accurately because the mass of this nucleus
is very large in comparison with the light particle y. Therefore, by using the law of conservation of linear
momentum we can calculate the Q-value.
Let the recoiling nucleus NR be emitted with a velocity VR at an angle aR and the particle y is emitted with
the velocity vy at an angle ay , as shown in Fig. 13.7. Now, from the law of conservation of momentum we
get the equations
MRVR cos aR + myvy cos ay = mbvb
or MRVR cos aR = mbvb – myvy cos ay (vii)
and myvy sin ay – MRVR sin aR = 0
or MRVR sin aR = myvy sin ay (viii)
By squaring and adding Eqs. (vii) and (viii), we get

MR2VR2 = mb2vb2 + my2vy2 – 2mbmyvbvy cos ay (ix)
1 1 1
By using the kinetic energy relations U b = mb vb2 , U y = m y v 2y and U R = M RVR2 in Eq. (ix), the expression
for UR becomes 2 2 2
mbU b m yU y 2
UR = + - (mb m yU bU y )1/2 cos a y (x)
MR MR MR
Now, substituting this relation into Eq. (v), we obtain for Q-value
Ê my ˆ Ê m ˆ 2
Q = U y Á1 + - U b Á1 - b ˜ - (mb m yU bU y )1/2 cos a y (xi)
Ë M R ˜¯ Ë MR ¯ MR
13.10.3 threshold energy
The above relation shows that the Q-value is independent of the mass of the target nucleus MT and the
kinetic energy of the recoiling nucleus, UR. As mentioned earlier, Q is positive for an exoergic reaction and
from Eq. (v) the final kinetic energy is the sum of Q and Ub. If Ub = 0, UR + Uy = Q > 0. Hence, we can say
that an exoergic reaction is energetically possible if the bombarding particle has zero kinetic energy. On the
other hand, endoergic reaction is energetically possible only when Ub > |Q|. The minimum kinetic energy of
the bombarding particle which is necessary to initiate the endoergic reaction is called threshold energy. We
can calculate this energy by using the centre-of-mass coordinate system, in which the linear momentum is
always zero before and after the reaction. If U¢b be the kinetic energy of the incident particle in this coordinate
system, the endoergic reaction is energetically possible only if
U¢b ≥ |Q| (xii)
mb M T
In terms of the reduced mass mr of the incident particle and the target nucleus, given by mr = , the
above condition can be written as mb + M T
1 mb M T 2
vb ≥ |Q| (xiii)
2 mb + M T
1 Ê M + mb ˆ
or mb vb2 ≥ Á T |Q| (xiv)
2 Ë M T ˜¯
Hence, the minimum energy required for endoergic reaction to take place, i.e., threshold energy, should be
Ê M + mb ˆ Ê m ˆ
(U b ) min = Á T ˜ |Q| = Á1 + b ˜ |Q|
Ë MT ¯ Ë MT ¯
This is clear from the above relation that the threshold energy is greater than the Q-value or the disintegration
È m ˘
energy by a factor of Í1 + b ˙ .
Î MT ˚
13.11 NucLear fissiON LO9

The phenomenon of breaking of a heavy nucleus into two or more light nuclei of almost equal masses
together with the release of a huge amount of energy is known as nuclear fission. The process of nuclear
fission was first discovered by the German scientists, Otto Hahn and Strassman, in the year 1939. In this
process, when uranium nucleus (23592U) was bombarded with slow neutrons, this nucleus was found to split
up into two radioactive nuclei which were identified as isotopes of barium (14456Ba) and krypton (8936Kr). It is
given by the following nuclear reaction.
235
92 U + 10 n Æ ÈÎ 92
236 ˘
U ˚ Æ 144
56 Ba +
89
36 Kr + 310 n + energy
It is not that barium and krypton are the only isotopes to be obtained by the fission of 235U. Actually, this is a
very complicated phenomenon and more than 100 isotopes of over 20 different elements have been obtained
in it. All these elements fall in the middle 75 to 160 mass number region of the periodic table.
13.11.1 Theory: Liquid-drop Model

The mechanism of nuclear fission was first explained by Bohr and Wheeler on the basis of liquid-drop
model of the nucleus. According to this model, the nucleus is assumed to be similar to a liquid drop, which
remains in equilibrium by a balance between the short-range, attractive forces between the nucleons and the
repulsive electrostatic forces between the protons. This inter-nucleon force gives rise to surface tension forces
to maintain a spherical shape of the nucleus. Thus, there is a similarity in the forces acting on the nucleus and
liquid drop. When nucleus-drop captures slow or thermal neutron, oscillations set up within the drop. These
oscillations tend to distort the spherical shape so that the drop becomes ellipsoid in shape (Fig. 13.11). The
surface-tension forces try to make the drop return to its original spherical shape while the excitation energy
tends to distort the shape still further. If the excitation energy and hence oscillations are sufficiently large, the
drop attains the dumbbell shape (Fig. 13.11). The Coulombic repulsive forces then push the nucleus into two
similar drops. Then each drop (bell) tries to attain the shape for which the potential energy is minimum, for
example, spherical shape.
144
Ba
56
n
n n
n
235
U
92
236
U Dumbbell 89
92 Kr
Shape 36
Figure 13.11
13.11.2 Nuclear energy
In nuclear fission, a huge amount of energy is liberated, which is known as nuclear energy. An estimation of
this energy can be made as follows. The mass of
235
92U + 01n = 234.99 amu + 1.01 amu = 236.0 amu
Similarly, the mass of
144 89
56 Ba + 36 Kr + 3 10 n = 143.87 amu + 88.90 amu + 3 ¥ 1.01 amu
= 235.80 amu
\ Mass defect Dm = 236.00 – 235.80 = 0.20 amu
According to Einstein’s mass-energy relation, E = mc2, 1 amu mass is equivalent to 931 MeV energy. So
energy released in each fission process = 0.20 ¥ 931 ª 190 MeV. This energy is millions of times more than
what we get by any chemical reaction.
13.11.3 chain reaction
When uranium is bombarded by neutrons, each uranium nucleus is broken into two nearly equal fragments
and a huge amount of energy is liberated and two or three fresh neutrons are emitted. If the conditions
are favourable, these neutrons take part in the fission of other uranium nuclei in the same way. This leads
to a chain of nuclear fissions which continues till the whole of uranium is fissioned within a fraction of
time (Fig. 13.12). Thus the energy produced in nuclear fission goes on multiplying. This energy takes a
tremendous magnitude very soon and is released as a violent explosion. Such a chain reaction is known as
uncontrolled chain reaction. This happens in an atom bomb.
If by some means, the reaction is controlled in such Ba

a way that only one of the neutrons produced in each
fission is able to cause further fission, then the fission U
Ba
process is slow and the energy is released steadily. This
chain reaction is known as controlled chain reaction. Slow Kr
This happens in nuclear reactors. U
Neutron Ba
13.11.4 critical size of Nucleus
Kr
U
We consider a uranium in sphere shape of radius r,
which has N1 as the number of neutrons produced in a
Kr
given time interval, N2 as the number of neutrons lost in
non-fission process and N3 as the number of neutrons Figure 13.12
escaped through the surface in the same time interval. N1
and N2 will be proportional to the volume, whereas N3 will be proportional to the surface area of the sphere,
i.e.,
4
N1 µ p r 3 = k1r 3
3
4
N 2 µ p r 3 = k2 r 3
3
and N3 µ 4p r 2 = k3 r 2
where k1, k2 and k3 are proportionality constants. The chain reaction is possible only when the rate of emission
of neutrons is greater than the total number of neutrons absorbed within the substance and going out of the
substance, i.e.,
N1 > N2 + N3
fi k1r3 > k2r3 + k3r2
or (k1 – k2)r > k3
k3
or r> = k (say)
k1 - k2
Reactor Pressure Vessel
where k is known as the critical size of the
Control Rods
nucleus. Thus, in order to achieve a self- Electric
Generator
sustained chain reaction, the size of the sample
must be greater than a critical value k. Below Turbine
this critical value the chain reaction will stop.
Cooling Water
13.11.5 Nuclear reactor
It is a device that produces a self-sustained
and controlled chain reaction in a fissionable Pump Pump
Fuel Rods
material. One type of nuclear reactor is Steam Generator
shown in Fig. 13.13. A modern reactor has
the following important parts. Figure 13.13
(i) Fuel The fuel plays the key role in the operation of the reactor. The fissionable material is known
as fuel. Generally, 235U and 239Pu can be used as fuel.
(ii) Moderator It is used to slow down the neutrons to thermal energies by elastic collisions between its nuclei
and the fission neutrons. Heavy water, graphite or beryllium oxide are commonly used for this purpose.
Heavy water is the most suitable moderator.
(iii) Control Rods To control the fission rate in the reactor, we use cadmium and boron rods. Cadmium
and boron are good absorbers of slow neutrons. These rods are fixed in the reactor-walls. When they
are pushed into the reactor, the fission rate decreases and when they are pulled out the fission rate
gets increased.
(iv) Shield The various types of intense rays, like a-, b-, g-rays in radioactivity are emitted from the
reactor. These rays may be injurious to the health of people working near the reactor. For protection,
the reactor is therefore surrounded by a concrete wall of about 2 meter thick and containing high
protection elements like iron.
(v) Coolant The reactor generates heat energy due to the fission reaction which is removed by means
of a cooling agent. For this purpose, air, water, carbon dioxide etc. are generally used as coolant.
Coolant is circulated through the interior of a reactor by a pumping system.
(vi) Safety Device If the reactor begins to go too fast, a special set of control rods, known as shut-
off rods drop inside automatically. They absorb all the neutrons so that the chain reaction stops
immediately.
13.11.5.1 Working of Nuclear Reactor

To start the reactor, no external source is required. Even a single neutron is capable of starting fission,
although few neutrons are always present there. The reactor is started by pulling out the control rods. Then
the neutron strikes 235U nucleus and fission it along with the emission of two or three fast neutrons. These
neutrons are slowed down by moderator (graphite), after which they induce further fission of 235U. The
reaction once started is controlled with the help of control rods by moving them inside and outside.
13.11.5.2 Applications of Nuclear Reactor

The nuclear reactors are used mainly for the following purposes.
(i) Generation of energy
(ii) Production of 239Pu
(iii) Production of neutron beam
(iv) Production of radioisotopes
13.12 NucLear fusiON

Nuclear fusion is nothing but the formation of a heavier nuclide by the fusing of two light nuclides. In this
process, the mass of the product nuclide is generally less than the sum of masses of the nuclides which
are fused. Therefore, as per Einstein’s mass energy relation E = mc2, an enormous amount of energy is
released, which is called nuclear energy. The first artificial fusion reaction was the hydrogen bomb which
was tested in November 1952. Fusion reactions are thermonuclear reactions which occur at extremely high
temperatures. For example, in order to fuse deuterium (21H) and tritium (31H), the force of repulsion (called
Coulomb potential barrier) of these two positively charged particles must be overcome.
The following fusion reaction is possible for the fusion of two heavy hydrogen nuclides 21H.
2
1H + 21H Æ 31H + 11H + 4.0 MeV (energy)
The nucleus of tritium (31H) can again fuse with heavy hydrogen nucleus
3
1H + 21H Æ 42H + 10n + 17.6 MeV (energy)
Thus, the combined form is
2
1H + 21H + 21H Æ 42He + 10n + 11H + 21.6 MeV (energy)
From the above equation, it is clear that three deuterium nuclei fuse together to form a helium nucleus and
liberate 21.60 MeV energy which is obtained in the form of kinetic energy of proton (11H) and neutron (10n).
The above reaction can also be possible in the following way:
2
1H + 21H Æ 32H + 10n + 3.30 MeV (energy)
3
2H + 21H Æ 42He + 11H + 18.30 MeV (energy)
13.13 cONtrOLLed fusiON LO9

Nuclear fusion is the main energy source that powers the sun and stars. As discussed earlier, in fusion the
nuclei of light elements (for example, hydrogen) fuse together to make heavier elements (for example, helium),
resulting in an enormous amount of energy. All over the world, efforts have been made to achieve this fusion
in a controlled manner for the utilisation of energy. The International Thermonuclear Experimental Reactor
(ITER) is a new attempt that will use a tokamak concept. Tokamak is a doughnut-shaped vessel in which a
strong, helical magnetic field guides the charged particles. Hence, a magnetic configuration is used to create
and maintain the conditions for controlled fusion reactions on earth. In ITER, superconducting magnet coils
around a toroidal vessel will confine and control a mixture of charged particles (called plasma) and induce an
electrical current through it. In this configuration, fusion reactions will occur when the plasma is hot enough,
dense enough and contained long enough (Lawson criterion, discussed later) for the atomic nuclei in the
plasma to start fusing together. The idea behind controlled fusion is to use magnetic fields to confine a high
temperature plasma of deuterium and tritium.
13.13.1 Plasma: the fourth state of Matter

You are aware of three states of the matter. These are solid, liquid and gas. In solids, the atoms are packed very
close to each other and are fixed at definite positions. These are connected with each other by the interatomic
forces. When we supply energy to the atoms of the solids, they start oscillating about their equilibrium
positions and as a result, the interatomic forces become weaker and the atoms are separated significantly. This
way the solid takes the form of liquid. The liquid has a specific volume but does not have a specific shape. So
it changes shape according to the shape of the container in which it is kept. Now if we again supply energy
to the atoms, the interatomic forces become insignificant, the atoms get separated and start moving freely. It
means the liquid has taken the form of gas. In gas, the atoms are not connected with each other and hence can
move in any direction. The gas does not have specific shape and volume also. It takes the shape and volume
of the container in which it is kept. If more energy is supplied to the gas species (atoms or molecules), the
electrons from the outermost level of the atoms get easily detached and hence the atoms become ionised. As
a result, we are left with the collection of ions, electrons and some neutrals (atoms that are not ionised). This
collection of charged and neutral particles is referred to as plasma, which is sometimes called the fourth state
of matter. This is because it is found in natural conditions; for example, the gases near the sun are always in
ionised state that qualify for plasma.
The species of the plasma being charged are connected with each other by the electromagnetic forces. This
can be explained as follows. Since the charges separated with each other give electric field, the plasma
species produce electric field. However, the separation of charges of plasma is not fixed (as the species do
not remain stationary; they keep on moving/oscillating). So, this electric field is a time-varying field, which

   ∂E
will generate magnetic field according to the Maxwell’s fourth equation — ¥ H = J + e 0 . However,
∂t
the motion of charges generates current and hence the magnetic field. In view of this, the plasma species
produce a time-varying magnetic field, which will induce electric field according to the Maxwell’s third

  ∂B
equation — ¥ E = - . Therefore, it can be said that the plasma species are connected with each other by
∂t
the electromagnetic fields. Since the number of ions and electrons in the plasma is almost equal, the plasma
as a whole is neutral. Since we cannot neglect the internal forces at the same time, the plasma is however
quasineutral. Moreover, if we attempt to perturb a part of the plasma, the whole body of the plasma will
get perturbed due to the connection of all the species with each other. This property is known as collective
behaviour of the plasma. Therefore, an ionised gas can qualify for plasma state, if it is quasineutral and shows
collective behaviour.
Another interesting property of the plasma is its ability to shield out the field which is applied on it. This
happens when we insert the electrodes of a battery into the plasma. The positive (negative) electrode
attracts the electrons (ions) whose number is decided by the charge carried by the electrode. Under this
situation, an electron cloud is developed around the electrode which shields/cancels the external field.
The thickness of this electron cloud is known as Debye length lDe. Since electrons are light species com-
pared with ions, the shielding is generally done by the electrons only. Clearly the field exists within the
cloud or the Debye sphere (sphere with the radius lDe). Now imagine if the Debye length is much less
than the dimension (L) of the plasma, the bulk of the plasma will remain neutral. Therefore, the required
condition for quasineutrality is lDe << L. Moreover, if the number of electrons in the Debye sphere (say,
NDe) is much larger than unity, i.e., NDe >> 1, the condition of collective behaviour will be fulfilled.
Any distance in the plasma system is measured in terms of Debye length lDe and the time is measured in
terms of inverse of plasma frequency fpe. The plasma frequency is nothing but the natural frequency of
the plasma, the same as all the materials have their natural frequencies. Actually, this is the frequency of
oscillations made by the electrons about their equilibrium positions. The Debye length lDe and the plasma
frequency fpe in SI system of units are given by
Ê e kT ˆ 1 Ê n0 e 2 ˆ
l De = Á 0 2e ˜ and f pe =
Ë n0 e ¯ 2p ÁË e 0 me ˜¯
Here k is the Boltzmann constant (= 1.38¥10–23 J/K), n0 is the plasma density, which is the common density
of ions (ni) and electrons (ne), i.e., n0 = ni = ne, Te is the electron temperature, e is the electron charge and me
is its mass.
In plasmas, generally we do not talk specifically about the temperature of the ions and electrons, but we
focus on their energies. That is, the temperature is written in terms of energy. For example, 1 eV energy of the
electron would be equal to its thermal energy kTe (for a 2-D system). So
1 eV = kTe
or 1.6 ¥ 10–19 (J) = 1.38 ¥ 10–23 ¥ Te (J/K)
or Te = 11,600 K
It means 1 eV of energy is equivalent to 11,600 K temperature. In laboratory plasmas, generally the electron
temperature varies from 1 eV to 5 eV. For a plasma with density 1018 /m3 and temperature 2 eV, the Debye
length is of the order of mm and the plasma frequency is of the order of GHz (109 Hz).
13.13.2 ignition temperature
A hot plasma at thermonuclear temperature loses a considerable amount of its energy in the form of radiation.
Therefore, it is required that the nuclear fusion produce more energy than is lost from the radiation of the
plasma. This requirement determines the minimum temperature for a nuclear fusion reactor to be self-
sustained. This temperature is called ignition temperature at which alpha-particle heating can sustain the
fusion reaction. As the temperature is increased, the production energy as well as the radiation losses get
increased. However, the fusion-energy production increases faster than the radiation loss.
13.13.3 Lawson criterion
For obtaining a net yield of energy from a fusion reaction, it is required, that in addition, to providing a
sufficiently high temperature to enable the particles to overcome the Coulomb barrier, this temperature must
be maintained for a sufficient confinement time and with a sufficient ion density. The overall conditions that
must be met for a yield of more energy than is required for heating of the plasma are generally stated in terms
of the product of ion density (n0) and confinement time (t). This condition is called Lawson criterion. For
deuterium-tritium (DT) fusion, it is
n0t ≥ 1014 sec/cm3
However, for deuterium-deuterium (DD) fusion, the Lawson criterion reads
n0t ≥ 1016 sec/cm3
13.13.4 Fusion by Inertial Confinement

Inertial confinement fusion (ICF) is an attempt for generating commercial energy using controlled
thermonuclear explosions. In this scheme, heavy isotopes of hydrogen, called fuel, are heated to temperatures
of around 10 keV. Actually, laser or particle beams are focused onto the surface of a capsule (diameter
of few millimetres), containing a small quantity of fuel. Due to evaporation and ionisation of the outer
layer of the material, a plasma crown is formed, which expands and, as in a rocket, generates an inward
moving compression front which heats up the inner layers of the material. This way, the core of the fuel is
compressed to as much as one thousand times its liquid density. Then ignition takes place when the core
temperature reaches about one hundred million degrees. Thermonuclear combustion spreads rapidly through
the compressed fuel and produces energy equivalent to several times the amount deposited on the capsule by
the laser or particle beams. The period of time during which these thermonuclear reactions occur is limited by
the inertia of the fuel itself. Because of this reason, such type of fusion is called fusion by inertial confinement
or inertial confinement fusion.
There are the following four basic phases of an inertial fusion pellet implosion.
(i) Irradiation In this process, irradiation of the pellet surface and formation of plasma is achieved with
the help of an intense energy beam such as laser. For achieving this, the laser is bombarded upon the pellet
(Fig. 13.14a).
(ii) Compression The compression of the pellet and fuel is driven by rocketlike blowoff of the surface material
(Fig. 13.14b).
(iii) Ignition The central fuel core is ignited to about 1000 to 10,000 times liquid deuterium-tritium
density and a temperature of 108°C (Fig. 13.14c).
(iv) Burning While the compressed fuel is inertially confined, it is burnt to achieve the fusion (Fig. 13.14d).
Figure 13.14
13.13.4.1 Laser Fusion
Laser fusion works on the concept of inertial confinement. If technically feasible, it would eliminate the
problems of magnetic instabilities. There are following two ways to achieve laser fusion.
(i) Laser-Gas-Fusion In this mechanism, CO2 laser beam is used to ionise and heat a long column of
gaseous deuterium and tritium at the density of n0 = 1017 /cm3. Here, the light of the laser is absorbed by
a process known as inverse bremstrahlung. This is because of the resistive damping of light wave due to
electron ion collisions. However, this process is not sufficient as for the density n0<< nc, where nc is the
critical density, the absorption length is very large (in kilometres).
(ii) Laser-Pellet-Fusion In this mechanism, laser light is focused on to a small pellet of solid deuterium–
tritium (DT), which has a number density n0 ª 5 ¥ 1022/cm3 and mass density r = 0.2 g/cm3. In this case since
n0 >> nc, the radiation is reflected as soon as plasma density of 1021/cm3 is formed on the pellet surface. This
depends on anomalous absorption by parametric decay instability to ionise the rest of the pellet and heat it to
a high temperature of 10 keV. As mentioned earlier, at this energy nuclei can penetrate the Coulomb potential
barrier and hence the nuclear reactions can take place.
13.13.5 Magnetic confinement
The basic problem in achieving controlled fusion is to generate plasma at very high temperatures and hold its
species (particles) together long enough for a substantial number of fusion reactions to occur. Since the
plasma is a mixture of charged particles, it can be

controlled and influenced by external magnetic fields.
Charged particles gyrate around the magnetic field
lines and also move along these field lines, as shown
in Fig. 13.15. These particles move in such a way that
the direction of magnetic field generated by their
motion is opposite to the direction of the applied field.
It is well established that the frequency at which a
charged particle, of charge q and mass m, gyrates in
the presence of magnetic field B0 is given by
| q | B0
wc =
m
If v^ be the perpendicular component of its velocity
with the direction of the magnetic field, the radius Figure 13.15
of the circle (in which the particle gyrates), called
Larmour radius, can be given as
v^ v m
rL = = ^
w c |q| B0
It means depending upon the strength and direction of the magnetic field, we can confine these charged
particles. Thus, these particles are confined by the application of magnetic field. However, if these field lines
are open–ended, theses particles are lost. In order to stop this loss of the particles at both ends, the magnetic
field lines can be closed to a ring. This configuration, which can be established by arranging a set of coils in
a ring, is called a torus.
13.14 ParticLe acceLeratOrs LO10

As discussed, in inertial confinement fusion, laser or particle beams are focused onto the surface of a pellet,
containing a small quantity of fuel, in order to evaporate and ionise the outer layer of the material. So highly
accelerated particles will contribute to achieve the controlled fusion. In this context, the acceleration of
particles and hence the particle accelerators are desirable. A particle accelerator is a device where we use
electric fields to propel electrically charged particles to high speeds and to contain them in well-defined
beams. An ordinary cathode ray tube (CRT) television set is a simple form of an accelerator. Basically, there
are two types of accelerators, namely, linear accelerators and circular accelerators.
13.14.1 Linear accelerator
In a linear accelerator, abbreviated as LINAC, charged particles are accelerated in a straight line with a target
of interest at one end. Linear accelerators operate on the general principle as used in a Van deGraaff generator
except that now a particle is exposed to a series of electrical fields, each of which increases the velocity of
the particle.
Typically, a linear accelerator consists of a few hundred or a few thousand cylindrical metal tubes which
are arranged one in front of another (Fig. 13.16). These tubes are electrically charged such that each carries
a charge opposite to that of the tube on either side of it. For example, tubes 1, 3, 5, etc. might be charged
positively and tubes 2, 4, 6, etc. charged negatively. Now imagine that an electron, which is negatively
charged, is introduced into a linear accelerator just in front of the first tube. Under the said configuration, the
electron is attracted by the first tube and is accelerated toward it. Then the electron passes into that tube. Once
inside the tube, the electron no longer feels any force of attraction or repulsion. So it merely drifts through the
tube until it reaches the opposite end. Because of this behavior, the cylindrical tubes in a linear accelerator
are generally referred to as drift tubes.
5
4
3
2
1
RF Oscillator
Ion Source Drift Tubes Vacuum Chamber
Figure 13.16
If the electron after leaving the first tube sees the next tube as positively charged, it will further accelerate. In
view of this, the moment that the electron leaves the first drift tube, the charge on all drift tubes is reversed.
So, tubes 1, 3, 5, etc., are now negatively charged and tubes 2, 4, 6, etc., are positively charged. Therefore,
the electron exiting the first tube now finds itself repelled by the tube it has just left. At the same time, it feels
attracted by the second tube. These forces of attraction and repulsion provide a kind of kick that accelerates
the electron in a forward direction in steps through different tubes.
As the electron moves through the linear accelerator, the electric charge on all drift tubes reverses in a regular
pattern. As mentioned earlier, the electron is repelled by the tube behind it and is attracted to the tube ahead
of it. This way the electron gains energy at every step. As a result, the electron moves faster in each new tube
it enters. Hence, it will cover a greater distance in the same amount of time. In order to make sure that the
electron exits a tube at just the right moment, each tube is made slightly longer than the one before it.
Stanford Linear Accelerator, located at the Stanford Linear Accelerator Center (SLAC) in Stanford, California,
is the largest LINAC in the world. This accelerator is 3 kilometres in length that holds 82,650 drift tubes along
with the magnetic, electrical, and auxiliary equipment needed for the machine’s operation. In this accelerator,
the electrons have been found to be accelerated up to 32 GeV (32 ¥ 109 eV).
13.14.2 cyclotron
The cyclotron is a particle accelerator which is also known as Lawrence Cyclotron, as it was conceived by
Lawrence in 1929. A cyclotron consists of two large dipole magnets designed to produce a semicircular
region of uniform magnetic field, pointing uniformly downward. Because of their D-shape these are called
D’s. The two D’s are placed back-to-back with their straight sides parallel but slightly separated, as shown
in Fig. 13.17.
Now in order to produce an electric field across this gap, we apply an oscillating voltage. Particles, which are
injected into the magnetic field region of a D, trace out a semicircular path until they reach the gap. However,
as the particles pass across the gap they are accelerated by the applied electric field. After gaining energy,
these particles follow a semicircular path in the next D with larger radius. Then they reach the gap again but
with the opposite direction of motion. Hence, the

electric field frequency must be such that the
direction of the field is reversed by the time of
arrival of the particles at the gap. This way the
particles are accelerated by the field in the gap. After
gaining energy, these particles enter the first D
again. Because of larger speed they trace a larger arc
and so they always take the same time to reach the
gap. This way a constant frequency electric field
oscillation continues to always accelerate the
particles across the gap. However, there is a
limitation on the energy that can be achieved in such
a device. This depends on the size of the magnets
Figure 13.17
that form the D’s and the strength of their magnetic
fields.
The cyclotron uses electric and magnetic fields and the whole accelerator remains in a uniform magnetic
field. In this accelerator, charged particles moving in the field feel a force acting at 90o to their direction of
motion. Hence, they move in circles. Here the Lorentz force due to the magnetic field provides the necessary
centripetal force for the circular motion with radius R. It means for a charged particle of charge q and mass
m circulating with velocity v,
mv 2 Ê mˆ
qvB = or R=Á ˜v
R Ë qB ¯
This equation says that the particles move in circles with radii proportional to their speeds.
Hence the time period of this motion
2p R 2p Ê m ˆ 2p m
T= = Á ˜ v=
v v Ë qB ¯ qB
Inverse of this time period gives the cyclotron frequency as
qB
f =
2p m
Clearly, it is not the function of the velocity v. This expression also shows that similar particles of same charge
and masses will all orbit at the same rate, regardless of their speed.
The particles’ orbits inside two metal semicircles are called Dees. These are connected to opposite terminals
of an ac supply which is set to the cyclotron frequency calculated above. The particles are accelerated in
bunches by the electric field between the plates. Since the cyclotron frequency is same for all the particles,
they all will be accelerated by the field. It is clear from R = mv/qB that the accelerated particles will move in
larger and larger circles until reaching a deflector at the edge of the machine which directs them to the target.
qBRmax
The maximum speed of the particles corresponding to maximum radius of circle would be vmax = .
m
1 2 q 2 B 2 Rmax
2
Therefore, maximum energy of the particles would be mvmax = . It is clear from this expression
2 2m
that for obtaining high energy, the strength of the magnetic field should be large and the radius of machine
also should be as large as possible. Since all the particles must orbit at the same frequency, whatever their
speed, the basic design of this accelerator is the real limiting factor. It is seen that as particles approach the
speed of light, they behave as if their mass is increasing. Therefore, their frequency becomes less and hence
they start to lag behind the oscillating electric field. So far, maximum energy gain achieved using a cyclotron
is about 20 MeV. A particle becomes relativistic once its kinetic energy is comparable to its rest energy. Since
the rest energy of the electron is only 500 keV, they cannot be accelerated to a useful energy in a cyclotron.
Lawrence and Livingston built the first cyclotron in 1932, which was about 0.3 m across in a magnetic field
of about 0.5 T. They could accelerate protons to roughly 1.2 MeV.
13.14.3 Betatron
The betatron can be thought of as a transformer with a ring of Electron Gun
electrons as the secondary coil (Fig. 13.18). In this accelerator,
Orbit
the magnetic field used to make the electrons move in a circle is
also the one used to accelerate them, as it is a rapidly alternating
field. However, the magnet must be carefully designed so that
Target
the field strength at the orbit radius (i.e., Borbit) is equal to half Doughnut
– B
the average field strength B linking the orbit, i.e., Borbit = .
2 β-Ray
We can change the flux by increasing the magnetic field. Since Beam
the flux links the loop of electrons, an induced emf accelerates Force from
the electrons. As the electrons get faster they need a larger Force from Guiding Field
Accelerating Field
magnetic field to keep moving at a constant radius. This is
provided by the increasing field. The field is changed by Figure 13.18
passing an alternating current through the primary coils. The particle acceleration takes place on the first
quarter of the voltage sine wave’s cycle. Although the last quarter of the cycle also has a changing field that
would accelerate the electrons, it is in the wrong direction for them to move in the correct circle. In order to
get effective acceleration, the target is bombarded with pulses of particles at the frequency of the ac supply.
The particles gain maximum energy when the magnetic field is at its strongest value. However, the formula
used for the cyclotron will not be appropriate for betatrons because the electron will be relativistic. If the total
energy is much greater than the rest energy, then E = pc is a good approximation. As the centripetal force is
again provided by the Lorentz force, we have
mv 2
= qvB
R
Here, B is the strength of the magnetic field which is needed to keep the particle in the orbit. From the
above relation, we obtain the maximum momentum as p = qBR. The maximum possible energy would be
E = pc = RqcB, where c is the speed of light.
The formula for the electron’s momentum can also be derived by using Faraday’s law of electromagnetic
df
induction. Since emf is equal to N , we can write
dt
dB
emf = NA
dt
 
2
For N = 1 and A = pR together with emf as  Ú E ◊ dl , we obtain the electric field E from the above relation
as below.
  dB
Ú E ◊ dl = p R2
dt
dB
or E ◊ 2p R = p R 2
dt
R dB
or E=
2 dt
The magnitude of force on the electron will be given by F = qE. Now as per Newton’s second law of motion,
we obtain
dp qR dB dp
F= or =
dt 2 dt dt
qRB
We can obtain the momentum of the particle if we integrate the above equation. This gives p = . Since
– 2
B = 2Borbit, the momentum p = qRBorbit.
Being compact in size, the betatrons are used in industry and medicine. Since cyclotrons cannot accelerate
electrons to useful energies, they are not as useful as betatrons are. A 315 MeV betatron was built in 1949 in
the University of Chicago.
13.14.4 Plasma–based Particle Accelerators

Plasma acceleration is a technique for accelerating charged particles like electrons, positrons, ions, etc.,
with the help of an electric field. This electric field is associated with an electron plasma wave, which is
produced either using electron pulses or by very short laser pulses. If we use laser pulses to excite this
wave, the technique is known as laser plasma acceleration. Through these techniques, we can achieve high-
performance particle accelerators along with much smaller size than conventional accelerators, for example,
RF linear accelerator (RF LINAC). This is because of the coherency, which is attained in such devices. These
devices show accelerating gradients of several orders of magnitude higher than the one of current particle
accelerators. For example, in an experimental laser plasma accelerator at Lawrence Berkeley National
Laboratory, the electrons can be accelerated to 1 GeV over about a 3.3 cm distance. However, the SLAC
conventional accelerator requires 64 m to reach the same energy. In another technique, called plasma wake
field acceleration (at SLAC) an energy gain of 42 GeV was achieved over 85 cm. It is believed that the plasma
acceleration technology could replace many of the traditional RF accelerators currently found in hospitals
and research facilities.
Since the plasmas are already in ionised state, unlike the conventional LINACs, there is no problem of
the electric field breakdown. The plasma can sustain the wave-breaking field E0 = cmewpe/e, where wpe
= 4p n0 e 2 /me is the electron plasma frequency (in CGS system of units) and n0 is the electron density.
For example, when n0 = 1018/cm3, the maximum attainable field E0 ª 100 GV/m. There are several variants
of plasma-based accelerators. Among those most investigated particle accelerators are plasma wake field
accelerators (PWFA), plasma beat wave accelerator (PBWA), laser wake field accelerator (LWFA) and self-
modulated laser wake field accelerator (SMLWFA).
13.14.4.1 Plasma Wake Field Accelerator (PWFA)
In order to understand the concept of wake field acceleration, it is necessary to understand what wake field is.
When a speed-boat travels in the water, it produces two types of waves, namely, bow waves and wake field
waves. The bow waves are conical waves having a tip at the front end of the boat. These are produced as the
velocity of the boat exceeds that of the water waves. The wake field waves are the waves set up at the back (or
wake) of the boat. These waves travel with the phase velocity equal to the velocity of the boat.
Figure 13.19
In a plasma wake field accelerator, the electron plasma wave is driven by one or more electron beams.
Effectively, the plasma wake fields can be excited by a relativistic electron beam (Fig. 13.19). This can be
achieved if the electron beam terminates in a time shorter than the plasma period wpe–1. In such a scheme,
the ratio of energy gain to the drive beam energy (called the transformer ratio Rt) is limited to Rt £ 2 for a
symmetric driving beam in the linear regime. However, it can be increased by using an asymmetric drive
beam. The idea of enhancing the wake field amplitude was also introduced by researchers, where they
proposed to use multiple electron drive bunches spaced at the plasma period.
13.14.4.2 Laser Beat Wave Accelerator (LBWA)
In the LBWA method, the plasma wave is exited by beating two optimal waves of significantly different
frequencies. Two long pulse laser beams having same direction of polarisation, and frequencies w1 and w2 are
used to resonantly excite a plasma wave (Fig. 13.20). When these pulses travel in a plasma of uniform density
n0(corresponding plasma frequency wpe), they will beat at Dw = w1 – w2 frequency. Here, the excitation
of plasma wave is done by appropriately adjusting the laser frequencies and plasma density such that the
resonance condition w1 – w2 = wpe is satisfied. Since the beat wave moves with laser pulse and plasma wave
also moves with a phase velocity equal to the group velocity of a laser wave, a properly placed bunch of
electrons with a velocity slightly lesser than the laser group velocity will get accelerated by the mechanism of
transfer of energy from wave to particle. However, in such mechanism, there is a problem of phase detuning
between the accelerated electrons and the plasma wave. This problem can be overcome if we make the use
of a transverse magnetic field.
LBWA
ω1 – ω2 = ωpe
Laser Pulse 1 Laser Pulse 2

(Frequency)
Frequency of Plasma Wave

Figure 13.20
13.14.4.3 Laser Wake Field Accelerator (LWFA)

In this technique, we use a short (time duration t £ 10–12 sec: 1 ps) laser pulse of very high intensity
(≥ 1018 W/cm2). As the laser pulse being an electromagnetic radiation propagates with the speed of light
c, its time duration t will correspond to the pulse length L = ct, which is nothing but the distance traveled
by the pulse in time t. As mentioned earlier, the wake is the electron plasma wave excited behind the pulse
and this wave propagates with the group velocity Laser Pulse
of the laser pulse. The group velocity of the pulse
is almost equal to the speed of light c. So, the Plasma Wave or Wake Wave
wavelength lp of this wake wave would be c/fpe,
where fpe is the frequency of the wave, equal to the LWFA
electron plasma frequency. Under this situation,
L = λp
if the laser pulse length L matches the plasma
wavelength c/fpe, then high amplitude wake field
will be produced due to resonance (Fig. 13.21). L
Therefore, a correctly placed traveling bunch of or λp
electrons can be accelerated by the longitudinal
Figure 13.21
field of the plasma wave or the wake field, which is
the field corresponding to the plasma wave.
13.14.4.4 Self–Phase Modulation LWFA
In this scheme, the electron plasma wave is excited resonantly by the modulation of the laser pulse envelope.
The self-modulated LWFA uses a single short (£ 1 ps) ultrahigh intensity (≥ 1018 W/cm2) laser pulse, as
in the case of LWFA. The self-modulated Single Pulse Modulated Pulse
LWFA, however, operates at higher plasma
densities than the standard LWFA. As we know
the plasma frequency fpe gets increased for the
higher density and hence the plasma wavelength Self–modulated
lp becomes smaller. So the condition L > lp, LWFA
which says that the laser pulse length is long
compared to the plasma wavelength, is satisfied L >��p
easily if the plasma density is higher. Moreover,
the laser power P is taken somewhat larger than P >�Pc
the critical power Pc. Under this situation, in this
high density regime, the laser pulse undergoes
a self modulation instability which causes the
pulse to become axially modulated at the plasma Plasma Period
period (Fig. 13.22). Hence, the plasma wave is Figure 13.22
generated resonantly by the modulated pulse.
s UmmarY
✦ The radius of the nucleus of an atom depends on the number of nucleons, A. R = r0A1/3, where
1.2 ¥ 10–3 £ r0 £ 1.48 ¥ 10–13 cm.
✦ The total angular momentum of the nucleons is given by I = L ± S; where I is a vector with magnitude
equal to maximum possible component of angular momentum in any direction, S is the total spin
angular momentum of all nucleons and L is the total orbital angular momentum.
✦ The nuclear magnetic moment mI is measured in terms of nuclear magneton mn, where value of
mN = 5.05 ¥ 10–27 J/wb/m2. –3mN £ mI £ +10 mN. mI = 0 for even-even- nuclei.
✦ The deviation of nucleus from its spherical symmetry results in electric quadrupole moment which is
responsible for the Electric property of a nucleus. The quadrupole moment Q of the nucleus is given
2
by Q = Z (b 2 - a 2), where Z is the magnitude of nuclear charge; 2b is the diameter along the axis of
5
symmetry and 2a is the diameter along the axis perpendicular to it.
✦ The statistical properties of assemblies of electrons, protons, neutrons, photons and atomic nuclei follow
the Quantum statistics. Systems of particles with anti-symmetric wave function such as electrons, protons,
neutrons and nuclei of odd mass number follow Fermi-Dirac statistics whereas systems of particles with
symmetric wave function such as nuclei with even mass number obey the Bose-Einstein statistics.
✦ The parity of a nucleus in a given state is related to the orbital angular momentum L. It is even for even
value of L and odd for odd value of L.
✦ Nuclear forces are of short range and are attractive in nature.
✦ Nuclear forces are charge independent. Extraordinarily stable nuclei are charge symmetric.
✦ The existence of mesons was predicted by Japanese physicist Yukawa. The mesons are massive particles
being constantly exchanged between two nucleons. Messons (p) can be neutral pions, negative pions
or positive pions.
✦ Binding energy is the energy released in the process of formation of the atom. It is the difference
between the sum of the mass of atom’s constituents and the mass of the atom.
✦ Binding energy is dependent on the atomic mass of atom. Atoms with every small atomic mass have
lesser Binding Energy per amu because majority of nucleons in this case are on surface, having lesser
binding energy. Atoms with large atomic mass also have less binding energy per amu since the binding
energy gets reduced by the Coulombic repulsion between protons.
✦ The nuclear stability is indicated by two factors; the N/Z ratio and the odd–even effect. N is the number
of neutrons and Z is the atomic number. The N/Z = 1 curve is called the stability curve and greater
deviation from the stability curve indicates an unstable nucleus. Even–even nuclei are the most stable,
followed by even–odd and odd–even. Odd–Odd nuclei are the least stable.
✦ Shell model and liquid drop models are the two most important models of nuclear structure. The nuclear
shell model considers the fact that there is a periodicity in the nuclear properties in terms of atomic
mass. There is a periodic increase in BE/A in case of nuclei with either Z or N equal to 2, 8, 20, 50,
82 and 126, called the magic numbers. The stability corresponding to the magic numbers is explained
by the formation of closed shells of protons or neutrons. The proton and neutron shells appear to be
independent of each other.
✦ In accordance with the shell theory of nuclear structure, the potential V(r) of nucleons is given by
Ê r2 ˆ
-Vo Á1 - 2 ˜ ; where Vo is the maximum potential between nucleons; r is the distance between the
Ë k ¯
nucleon and the centre of force and k is a constant.
✦ The shell model of nuclear structure has enabled the prediction of total angular momenta of nuclei,
occurrence of isomerism and zero quadrupole moment at proton number 2, 8, 20, 50, and 82.
✦ According to Liquid Drop Model, the nucleus is considered analogous to a drop of incompressible,
high density liquid. Taking into account the volume energy, surface energy, symmetry effect, Coulomb
energy and the odd–even effect of the nucleus; the Binding energy for a nucleus is given by—
0.584 Z ( Z - 1) 19.4 ( A - 2Z ) 2
BE = 14 A - 13.1 A2/3 - - + Ed MeV
A A
135 135
Ed = for even–even nuclei, - for odd–odd nuclei and O for even–odd and odd–even nuclei.
A A
✦ The discovery of neutron in 1932 solved a known puzzle related to the spin of the nitrogen-14 nucleus,
which was experimentally measured as 1 basic unit of angular momentum, but at that time, physicists
could not find any way to arrange 21 particles of 147N so as to give a spin of 1. However, the presence of
neutron as an uncharged particle in the nucleus with spin ½ solved this problem. Another importance of
neutron is in nuclear fission as a slow neutron initiates the fusion and the resulting reactions produce on
an average 2.4 neutrons. This way, a chain reaction is generated that allows a self-sustaining mode of
operation.
✦ Radioactivity is the disintegration of certain natural heavy elements, which is accompanied by the
emission of a-rays (positively charged helium nuclei), b-rays (fast electrons) and g-rays (short X-rays).
The ultimate end product of the radioactive disintegration process is an isotope of lead. The radioactivity
is of two types, namely, artificial radioactivity and induced radioactivity.
✦ If l be the disintegration constant or decay constant, N0 be the initial number of nuclides at time t = 0
and N be undecayed nuclei at time t, then the decay takes place as per the relation N = N0e–lt.
dN
✦ The activity A is the number of disintegrations per second of a sample. It is given by A = =
dt
lN0e–lt = lN.
✦ The half-life time, T1/2, of any sample is defined as the time interval in which the number of undecayed
0.693
atoms decreases by half. It is given by T1/2 = .
l
✦ The mean life time t of a nuclide is the reciprocal of its decay probability per unit time. The mean life
T
time and half-life time T1/2 are related to each other as t = 1/2 = 1.44 T1/2 .
0.693
✦ If a nucleus contains 210 or more nucleons, i.e., when nuclei are so large that the short-range nuclear
forces that hold them together are barely able to counterbalance the mutual repulsion of their protons, the
nuclei decay by the process of alpha (42He) decay in order to increase their stability. The a decay takes
place as per ZA X Æ ZA-- 42Y + 42 He relation, if a parent nucleus X disintegrates into daughter nucleus Y.
✦ Beta decay is a radioactive decay in which a beta particle that may be either an electron or a positron
is emitted. Actually, three mechanisms are involved in beta decay. These are b– decay (electron or b–
emission, or negatron emission), b+ decay (or positron emission) or electron capture, in which a nucleus
decays by capturing an extra nuclear atomic electron and the electron is disappeared because its mass
is converted into energy.
✦ By emitting alpha, beta or other particles, the nucleus disintegrates and is usually left in the excited
state. If the excited nucleus does not have sufficient energy to emit another particle, like an excited
atom, it returns to its ground state by emitting photons. These emitted photons from nuclei are called
gamma rays.
✦ The detection of nuclear radiations depends upon their interaction with matter and especially on
the excitation and ionisation processes. A nuclear-radiation detector is a device in which presence
of radiation induces physical change that is observable. These detectors include ionisation chamber,
Geiger–Mueller counter, scintillation counter, Wilson’s cloud chamber, bubble chamber and
semiconductor detector.
✦ Neutron absorption cross-section is the cross-section for a nuclear reaction which is initiated by
neutrons. The same way we can define neutron scattering cross-section. Neutron cross-section is the
measure of probability of scattering and absorption of neutron when it approaches a nucleus. This is
expressed in barns (1 barn = 10–28 m2).
✦ If a target material is bombarded by fast-moving particles such as protons, neutrons, electrons, deuterons
or alpha particles, then the target nuclei after the bombardment are usually different from what they
were before. The target nuclei may change their mass number or atomic number or both due to its
interaction with an incident particle. This is called a transmutation and the reaction is called nuclear
reaction.
The nuclear reactions are of two types: exoergic reaction and endoergic reaction. If the disintegration
energy or the Q-value is positive, then the nuclear reaction is called exoergic reaction. If Q-value is a
negative quantity, then the nuclear reaction is an endoergic reaction. The minimum kinetic energy of
the bombarding particle which is necessary to initiate the endoergic reaction is called threshold energy.
✦ The phenomenon of breaking of a heavy nucleus into two or more light nuclei of almost equal masses
together with the release of huge amount of energy is known as nuclear fission. The released energy in
this process is called nuclear energy. In order to achieve a self-sustained chain reaction, the size of the
sample must be greater than a critical value, which is called the critical size of the nucleus.
✦ Nuclear fission is achieved in a nuclear reactor, which produces a self-sustained and controlled chain
reaction in a fissionable material. The important parts of the reactor are fuel, moderator, control rods,
shield, coolant and safety device.
✦ Nuclear fusion is the formation of a heavier nuclide by the fusing of two light nuclides. The first
artificial fusion reaction was the hydrogen bomb which was tested in November 1952. Fusion reactions
are thermonuclear reactions which occur at extremely high temperatures. For example, in order to fuse
deuterium (21H) and tritium (31H), the force of repulsion (called Coulomb potential barrier) of these two
positively charged particles must be overcome.
✦ Nuclear fusion is the main energy source that powers the sun and stars. All over the world, efforts
have been made to achieve this fusion in controlled manner for the utilisation of energy. This is called
controlled fusion. International Thermonuclear Experimental Reactor (ITER) is a new attempt that will
use the concept of a tokamak, which is a doughnut-shaped vessel in which a strong, helical magnetic
field guides the charged particles. The idea behind controlled fusion is to use magnetic fields to confine
a high temperature plasma of deuterium and tritium.
✦ Plasma is a collection of charged and neutral particles. The charged particles are nothing but ions and
electrons. These ions and electrons are almost in equal numbers in the plasma and the plasma as a
whole is neutral. Since we cannot neglect the internal forces at the same time, the plasma is however
quasineutral. Moreover, if we attempt to perturb a part of the plasma, the whole body of the plasma
will get perturbed due to the connection of all the species with each other. This property is known
as collective behaviour of the plasma. Therefore, an ionised gas can qualify for plasma state, if it is
quasineutral and it shows collective behaviour.
✦ A hot plasma at thermonuclear temperature loses a considerable amount of its energy in the form of
radiation. In order to produce more nuclear fusion energy than is lost from the radiation, there is the
requirement of a minimum temperature for a nuclear reactor to be self-sustained. This temperature is
called ignition temperature at which alpha particle heating can sustain the fusion reaction.
✦ The overall conditions that must be met for a yield of more energy from the nuclear reactor than is
required for heating of the plasma are generally stated in terms of the product of ion density (n0) and
confinement time (t). This condition is called Lawson criterion. For deuterium–tritium (DT) fusion, it
is n0t ≥ 1014 sec/cm3. However, for deuterium–deuterium (DD) fusion, the Lawson criterion reads n0t
≥ 1016 sec/cm3.
✦ In inertial confinement fusion (ICF), heavy isotopes of hydrogen, called fuel, are heated to temperatures
of around 10 keV. Laser or particle beams are focused onto the surface of a capsule containing a small
quantity of fuel. Due to evaporation and ionisation of the outer layer of the material, a plasma crown
is formed, which expands and generates an inward moving compression front which heats up the
inner layers of the material. So the core of the fuel is compressed to as much as one thousand times its
liquid density. Then ignition takes place when the core temperature reaches about one hundred million
degrees. Thermonuclear combustion spreads rapidly through the compressed fuel and produces energy
equivalent to several times the amount deposited on the capsule by the laser or particle beams.
✦ Laser fusion works on the concept of inertial confinement. If technically feasible, it would eliminate the
problems of magnetic instabilities. There are two ways to achieve laser fusion, namely, laser-gas-fusion
and laser-pellet-fusion.
✦ The basic problem in achieving controlled fusion is to generate plasma at very high temperatures and
hold its species together long enough for a substantial number of fusion reactions to occur. Since the
plasma is a mixture of charged particles, it can be controlled and influenced by external magnetic
fields. Charged particles gyrate around the magnetic field lines and also move along these field lines;
so by applying the magnetic field with a proper configuration, we can confine the plasma. This is called
magnetic confinement.
✦ Highly accelerated particles contribute to achieve the controlled fusion, as the laser or particle beams
are focused onto the surface of a pellet, containing a small quantity of fuel, in order to evaporate and
ionise the outer layer of the material. Therefore, the acceleration of particles is a topic of interest. A
particle accelerator is a device where we use electric fields to propel electrically charged particles
to high speeds and to contain them in well-defined beams. An ordinary cathode ray tube (CRT)
television set is a simple form of accelerator. Various types of accelerators, namely, linear accelerator,
cyclotron, betatron, and plasma-based accelerators were discussed. Plasma-based accelerators include
plasma wake field accelerator, laser beat wave accelerator, laser wake field accelerator and self-phase
modulation laser wake field accelerator.
s olved e XamPles
E xamplE 1 In an absorption experiment with 1.14 MeV g-radiation from 65Zn it is found that 20 cm of
aluminium reduces the beam intensity to 3%. Calculate the half–value thickness and mass absorption
coefficient of Al for this radiation. Density of aluminium = 2700 kg/m3.
3
Solution Given I/I0 = 3% = , x = 20 cm = 0.20 m and r = 2700 kg/m3
100
The formula used is
I = I0 e–mx (i)
where m is linear absorption coefficient and x is thickness of aluminium.
I I
= e - m x or ln = - m x
I0 I0
I (ii)
or m x = ln 0
I
2 1
x = 0.2 m = m = m
10 5
Putting the value of x in Eq. (ii), we get
m Ê 100 ˆ
= ln Á = ln100 - ln 3
5 Ë 3 ˜¯
= 4.60517 – 1.0986
m
= 3.5065577
5
m = 17.5328 m–1
m 17.5328 2
Mass absorption coefficient = = m / kg
r 2700
or = 0.00649363 m2/kg
0.693
Half-value thickness ( x1/ 2 ) =
m
0.693
or x1/ 2 = = 0.0395 m
17.5328
E xamplE 2 In an absorption experiment with 1.1 MeV g-radiation from 65Zn it is found that 25 cm of Al
reduces the beam intensity to 2%. Calculate the half-thickness, and the mass attenuation coefficient of Al for
this radiation. Density of Al = 2700 kg/m3.
2 1 1
Solution Given I / I 0 = 2% = = , x = 0.25 m = m
100 50 4
Formula used is I = I0e–mx
or mx = ln ÊÁ I 0 ˆ˜
Ë I¯
m
= ln(I0/I) or m = 4 ln(50)
4
m = 4 ln 50 = 4 ¥ 3.912 m–1
= 15.648 m–1
Mass attenuation (or absorption) coefficient
m 15.648
= =
r 2700
= 5.8 ¥ 10 -3 m 2 / kg
0.693 0.693
and half-value thickness ( x1/ 2 ) = = = 0.0443 m
m 15.648
E xamplE 3 Half-life of 23
11 Na is 15 hours. How long does it take for 93.75 % of a sample of this isotope to
decay?
23
Solution Given Half-life of 11 Na, i.e., T1/2 = 15 hours.
0.693
The radioactive constant l =
T1/ 2
or 0.693 -1
l= hr
15
= 0.0462 hr -1
N
Now, = e - lt
N0
Here N0 is the number of atoms that existed in beginning and N is the number of atoms left behind after time t. Then
N 6.25
= = e - lt
N 0 100
1 Ê 1ˆ
or = e - l t or ln Á ˜ = - lt
16 Ë 16 ¯
ln(16) 2.7726
or t= =
l 0.04621
= 60 hrs
E xamplE 4 Half-life of a radioactive element is 4 years. After what time will the element present in specimen
reduce to 1/64 of its original mass?
Solution Given T1/2 = 4 yrs and N/N0 = 1/64

0.693 0.693
Decay constant (l ) = = = 0.17327 per yr
T1/ 2 4
N 1
and = = e - lt or lt = ln 64
N 0 64
ln 64
t= = 24.002326
0.17327
= 24 yrs
E xamplE 5 The half-life of a radioactive substance is 15 years. Calculate the period in which 2.5% of the
initial quantity will be left over.
2.5 1
Solution Given T1/2 = 15 yrs and N/N0 = =
100 40
0.693 0.693
Decay constant l = = = 0.0462 per year
T1/ 2 15
N 1
and = = e - lt or lt = ln(40)
N 0 40
ln 40 3.689
t= =
l 0.0462
= 79.85 yrs
E xamplE 6 How long does it take for 60% of a sample of radon to decay? T1/2 for radon is 3.8 days.
Solution Given T1/2 = 3.8 days. If 60% radon decays, it means 40% of it is left behind.
0.693
Decay constant (l ) =
3.8
= 0.18237
= 0.1824 d–1
Now N = N0e–lt
100
or 40 = 100 ¥ e - l t or lt = ln = ln 2.5
40
ln 2.5 0.9163
or t= =
l 0.1824
t = 5.024 d
E xamplE 7 Calculate the half-life time and mean life time of the radioactive substance whose decay constant
is 4.28 ¥ 10–4 per year.
Solution Given l = 4.28 ¥10–4 per year.
0.693
Half-life time T1/ 2 =
l
0.693
or T1/ 2 =
4.28 ¥ 10 -4
= 1619.16 yrs
1
Now mean life time t =
l
1
or t = = 2336.45 yrs
4.28 ¥ 10 -4
ExamplE 8 Find the half-life of a radioactive material if its activity drops to 1/64 of initial activity in 30 years.
1
Solution Given t = 30 years and A = A0 .
64
Activity A = A0e–lt
A0 1
= A0e - l t or = e - lt
64 64
or lt = ln 64
ln 64 4.1589
l= =
t 30
l = 0.1386 yrs -1
0.693 0.693
Half-life time T1/ 2 = =
l 0.1386
= 4.999 yrs
= 5.0 yrs
E xamplE 9 What is the decay constant of a nucleus whose half-life is 2.1 min?
Solution Given T1/2 = 2.1 min.
Decay constant
0.693 0.693
l= = per min
T1/ 2 2.1
= 0.33 min -1
E xamplE 10 Calculate the decay constant for 198Au whose half-life is 2.7 days. If at some time, a sample
contains 10 gm of 198Au, what would be its activity? Calculate decays occurring per second after 8 days.
–6
Solution Given T1/2 = 2.7 d, m = 10–6 g and t = 8 d.

Decay constant
0.693 0.693
l= =
T1/ 2 2.7 ¥ 24 ¥ 60 ¥ 60
= 2.971 ¥ 10 -6 sec -1 (i)
The number of nuclei in sample N can be calculated as
mass ¥ Avogadro number
N=
molar mass
(10-6 g )(6.02 ¥ 1023 atoms/mol)
=
(198 g/ mol)
= 3.04 ¥ 1015 atoms
But putting the value of l and N in A0 = lN we get the following for activity at t = 0
A0 = 2.971 ¥ 10–6 ¥ 3.04 ¥ 1015
= 9.032 ¥ 109 number of atoms disintegrated per second
Now Activity A = A0e–lt
A = (9.032 ¥ 109) ¥ e–2.054
= 9.032 ¥ 109 ¥ 0.12822
A = 1.1581 ¥ 109 decays/sec
E xamplE 11 What fraction of a sample is left after 3 half-lives?

Solution Let N0 be the original number of nuclei and N0/2 be the number of nuclei remains after half-lives.
N0 / 2 1
\ The fraction after half-lives = =
N0 2
N 0 /8 1
So fraction after 3 half-lives = =
N0 8
= 0.125
E xamplE 12 Ten milligrams of a radioactive substance of life period 2 years is kept for four years. How
much of the substance remained unchanged?
Solution The substance remains unchanged after 4 years would be
N 0 10 mg
N= = = 2.5 mg
4 4
E xamplE 13 One gram of radium is reduced by 2.1 mg in five years by a-decay. Calculate decay constant,
half-lives of sample, and average life.
Solution Given N0 =1.0 g and N = 1 – 0.0021 g = 0.9979 g.
N N0
Now = e - lt or lt = ln
N0 N
Ê 1.0 ˆ
ln Á
Ë 0.9979 ˜¯ ln (1.0021) 2.1022 ¥ 10-3
or decay constant l = = =
t 5 5
= 3.996 ¥ 10 -4 per year
= 0.0004 per year
0.693
Half-life time T1/ 2 = = 1732.5 yrs
0.0004
1 1
Average life t = =
l 4 ¥ 10 -4
= 2500 yrs
E xamplE 14 The activity of certain radio nuclide decreases to 15% of its original value in 10 days. Find its
half-life.
Solution Given N0 = 100, N = 15 and t =10 days
N
Now N = N 0e - l t or = e - lt
N0
ÊN ˆ Ê 100 ˆ
lt = ln Á 0 ˜ = ln Á
Ë N¯ Ë 15 ˜¯
1.897
l= = 0.1897 d -1
10
0.693 0.693
Half-life T1/2 = =
l 0.1897
= 3.65 d
E xamplE 15 What fraction of a radioactive isotope remains after 50 years, if its half-life is 13.3 years?
Solution Given t = 50 years and T1/2 =12.3 years.
Fraction of a radioactive isotope
N
= e - lt (i)
N0
where decay constant
0.693 0.693
l= =
T1/ 2 12.3 years
= 0.05634 yr -1
and lt = 2.817.
Now putting the value of lt in Eq. (i), we get

N
= e -2.817
N0
= 0.0598 = 0.06
E xamplE 16 Calculate the mass of 214Pb having radioactivity of 1 Curie. Half-life of 214Pb is equal to 26.8
minutes.
Solution Given T1/2 = 26.8 min = 26.8 ¥ 60 sec, and 1.0 Curie = 3.7 ¥ 1010 disintegrations/sec.
Let m gram mass of 214Pb has an activity of 1 Curie, then the number of atoms in m gram of 214Pb.
6.023 ¥ 1023 ¥ m
N=
214
0.693 0.693
Disintegration constant l = =
T1/ 2 1608 sec
and activity = 4.31 ¥ 10 -4 sec -1
A = Nl
6.023 ¥ 1023 ¥ m ¥ 4.31 ¥ 10 -4
So A=
214
6.023 ¥ 10 23
¥ m ¥ 4.31 ¥ 10 -4
3.7 ¥ 1010 =
214
-8
or m = 3.05 ¥ 10 g
E xamplE 17 Calculate the weight in grams of 214Pb from the half-life of 26.8 minutes when its activity is 106.
23
Solution As done in Ex. 16, N = 6.025 ¥ 10 ¥ m ,
214
0.693
decay constant l = = 4.31 ¥ 10 -4 sec -1
26.8 ¥ 60 sec
6.023 ¥ 1023 ¥ m ¥ 4.31 ¥ 10 -4
and activity A = Nl =
214
6.023 ¥ 1023 ¥ 4.31 ¥ 10 -4 ¥ m

or 106 =
214
or m = 8.24 ¥ 10 -13 g
E xamplE 18 One gram of 226Ra has an activity of one Curie. Calculate the mean life and half-life of radium.
Solution Given activity A = 1 Curie = 3.7 ¥ 1010 disintegrations per second,
6.023 ¥ 1023 ¥ 1
and N=
226
Activity A = Nl
A 3.7 ¥ 1010 ¥ 226
or l= =
N 6.023 ¥ 1023
l = 1.38 ¥ 10 -11 sec -1
1 1
Mean life = =
l 1.38 ¥ 10 -11
= 7.25 ¥ 1010 sec
0.693 0.693
and half-life T1/ 2 = =
l 1.38 ¥ 10 -11
= 5.02 ¥ 1010 sec
ExamplE 19 Calculate the activity of 0.1 mg sample of 90Sr at t = 9 sec when half-life period of 90Sr is 28 years.
Solution Given T1/2 = 28 yrs = 28 ¥ 365 ¥ 24 ¥ 60 ¥ 60 sec
= 8.83 ¥ 108 sec
Decay constant
0.693 0.693
l= =
T1/ 2 8.83 ¥ 108
= 7.85 ¥ 10 -10 sec -1
Number of disintegration per second = N0 – N
or N0 – N0e–lt = N0 – N0[1 – lt + …]
= N0[lt] = N0lt
6.023 ¥ 1023 ¥ 10 -4 È Avogadro numbers ¥ mass ˘

Now N0 = = 6.69 ¥ 1017 Í N 0 = ˙
90 Î Molar mass ˚
So number of disintegration per sec = N0 lt
= 6.69¥1017¥7.85¥10–10¥9
= 5.25 ¥ 108
E xamplE 20 The half-life of radium (226) is 1600 years and that of radon (222) is 3.8 days. Calculate the
mass of radon that will be in equilibrium with one g of radium.
Solution Half-life of radium (T1/2) = T1 = 1600 yrs and half-life of radon (T1/2) = T2 = 3.8 days.
Let N1 be the number of atoms in one g of radium, i.e.,
N
N1 = (i)
226
and suppose m is the mass of radon-222 which is in equilibrium with 1.0 g of radium. Then the number of atoms in
m gram radon, i.e.,
mN
N2 = (ii)
222
Again consider l1 and l2 as the radioactive decay constants for radium and radon, respectively, and T1 and T2 are
corresponding half-life periods. In equilibrium
0.693 0.693
N1l`1 = N 2 l2 or N1 = N2 (iii)
T1 T2
By using Eqs. (i) and (ii) in Eq. (iii), we get
N 1 Nm 1
=
226 T1 222 T2
1 1 m 1
or =
226 365 ¥ 1600 222 3.8
or m = 6.39 ¥ 10-6 g
E xamplE 21 How much energy would a g-ray photon have if it is to split an a-particle into a tritium 31H
1
and proton 1H?Given masses of 24He, 13H and 11H as 4.002603 a.m.u., 3.016056 a.m.u. and 1.007276 a.m.u.
respectively.
Solution Given ma = 4.002603 a.m.u., mt =3.016056 a.m.u. and mp = 1.007276 a.m.u.
According to the problem, the reaction may be
4
2 He + g = 13H + 11H (i)
By putting the values of masses of various constituents in Eq. (i), we get
4.002603 + g = 3.016056 + 1.007276
The mass of g-ray photon = 4.023332 – 4.002603
= 0.020729 a.m.u.
Equivalent energy of g-ray photon is = 0.020729 ¥ 931 MeV
= 19.298 MeV
E xamplE 22 A tritium gas target (13H) is bombarded with a beam of protons (11H) of kinetic energy 3MeV.
Determine Q value of the following reaction and specify the type of reaction.
1
1H + 13H Æ 23He + 01n + Q
Given m(11H) = 1.007276 a.m.u.; m(13H) = 3.016056 a.m.u.; m(01n) = 1.008665 a.m.u.; m(23He) = 3.016036 a.m.u.
Solution The given reaction is

1
1H + 13H Æ 23He + 01n + Q (i)
By putting the values of masses of various constituents in Eq. (i), we get
1.007276 + 3.016056 = 3.016036 + 1.008665 + Q
Q = 4.023332 – 4.024701
= –0.001369 a.m.u.
= –0.001369 ¥ 931.5 MeV
= –1.2745 MeV
The negative sign in Q value indicates that 1.2745 MeV is required for the reaction to occur. Hence, this reaction is
endoergic reaction.
ExamplE 23 Assuming that 200 MeV of energy is released per fission. Calculate the energy released in Joules
and also the heat produced by complete disintegration of 10 mg of235
92 U.
Solution The energy released per fission of235
92 U atoms
= 200 MeV = 200 ¥ 106 ¥ 1.6 ¥ 10–19 J

=3.2 ¥ 10–11 J
The number of atoms in one gram of235 23
92 U is given by Avogadro’s number, i.e., 6.023 ¥ 10 .
\ The energy released by one gram atom of the235
92 U
3.2 ¥ 10-11 ¥ 6.03 ¥ 1023

=
235
Hence the energy released by fission of 10 ¥ 10–3 g of235
92 U
3.2 ¥ 10 -11 ¥ 6.03 ¥ 1023 ¥ 10 -2
=
235
= 8.21 ¥ 108 J
8.21 ¥ 108
Heat produced = calories = 1.961 108 calories
4.186
E xamplE 24 Considering the average energy released per fission as 200 MeV, determine the energy released
by fission of 1.0 kg of 235U. Given Avogadro number as 6.03 ¥ 1026 per kg atom.
Solution The energy per fission of 235U atom is
= 200 MeV = 200 ¥ 106 ¥ 1.6 ¥ 10–19 J
= 3.2 ¥ 10–11 J
3.2 ¥ 10 -11 ¥ 6.03 ¥ 1026
The energy released by the fission of 1.0 kg of 235U is =
235
= 8.21 ¥ 1013 J
E xamplE 25 On an average, 1 GW electric power is required to enlighten a city. If a nuclear reactor of
efficiency 30% is used for the same purpose with 235U as a nuclear fuel, what amount of fuel would be
required per day. Consider the energy released per fission of 235U as 3.2 ¥ 10–11 J.
Solution Given energy required per second =1 GW
=1.0 ¥ 109 J,
Energy released per fission = 200 MeV
= 200 ¥ 106 ¥ 1.6 ¥ 10–19 J
= 3.2 ¥ 10–11 J
30
Efficiency of reactor is 30%, so the actual energy released per fission = 3.2 ¥ 10 -11 ¥
100
= 9.5 ¥ 10–12 J
1.0 ¥ 109
So, the number of fission required per second = = 1.04167 ¥ 1020
9.6 ¥ 10 -12
Number of atoms required undergoing fission per day =1.04167 ¥ 1020 ¥ 24 ¥ 60 ¥ 60
= 9.0 ¥ 1024
6.03 ¥ 1026
Number of atoms in one kg of 235 U =
235
= 2.5659574 ¥ 1024
Amount of fuel required per day for operation of reactor with 30% efficiency (weight of 235U for 9.0 ¥ 1024 atoms)
9.0 ¥ 1024
= = 3.51 kg
2.5659574 ¥ 1024
= 3.51 kg
E xamplE 26 In a nuclear reactor, the fission of 235U atom yields 200 MeV. If energy of 3.7 kg uranium is
consumed in a day, find the power output of the reactor assuming that the reactor is 20% efficient.
Solution The number of atoms in 1.0 kg of 235U
6.03 ¥ 1026
=
235
= 2.5659574 ¥ 1024
and the total number of atoms in 3.7 kg of 235U
= 3.7 ¥ 2.5659574 ¥ 1024
= 9.494 ¥ 1024
Energy released per fission = 200 MeV = 200 ¥ 106 eV
= 2.0 ¥ 108 ¥ 1.6 ¥ 10–19 J
= 3.2 ¥ 10–11 J
20
The efficiency of the reactor is 20%, so the net energy released per fission = 3.2 ¥ 10 -11 ¥
100
= 6.4 ¥ 10–12 J
The net energy released due to consumption of 3.7 kg of U235 per day = 9.494 ¥ 1024 ¥ 6.4 ¥ 10 -12 J/day
60.7616 ¥ 1012
= J/s
24 ¥ 60 ¥ 60
= 7.033 ¥ 108 W
= 0.703 GW
E xamplE 27 In a reaction, the energy is produced by the fusion of the three helium nuclei to form a 12
6 C
nucleus. How much energy is produced by each reaction? Consider the mass of helium atom, electron and
12
6 C as 4.00260 a.m.u., 0.00055 a.m.u. and 12.0000 a.m.u., respectively.
Solution Mass of 3 helium nuclei = 3 ¥ 4.00260 a.m.u.

= 12.00780 a.m.u.
Mass of 12
6 C
atom = 12.0000 a.m.u.
Mass defect (Dm) = 12.00780 – 12.0000
= 0.00780 a.m.u.
Energy produced = 0.00780 ¥ 931 MeV
= 7.2618 MeV
E xamplE 28 In an industry, the energy is produced using the fusion reaction 21H + 21H Æ 24H + energy. If the
efficiency of fusion reactor is 33%, calculate how much deuterium will be consumed per day for production
of 50 MW energy. Consider mass of 21H and 42He as 2.01478 and 4.00388 a.m.u. respectively.
Solution According to the fusion reaction, the mass difference
= 2.01478 + 2.01478 – 4.00388 = 0.02568 a.m.u.
\ Equivalent energy (i.e., energy produce per fission)
= 0.02568 ¥ 931 = 23.908 MeV
= 23.91 MeV
33
Efficiency of fusion reactor is 33% =
100
Energy output 33
i.e., =
Energy input 100
33
\ Energy output ¥ 23.91 = 7.89 MeV = 1.262 ¥ 10–12 J
100
The actual energy output per deuterium atom
1.262 ¥ 10 -12
= = 0.631 ¥ 10 -12 J
2
The number of deuterium atoms required per sec for production of 50 MW energy
50 ¥ 106 J/sec
= = 7.924 ¥ 1019 atoms/sec
0.631 ¥ 10 -12 J
Mass of one deuterium atom = 2.01478 a.m.u.
2.01478
= kg
6.03 ¥ 1026
= 3.3413 ¥ 10 -27 kg
The equivalent mass of deuterium atoms consumed in production of energy per second
= 7.924 ¥ 1019 ¥ 3.3413 ¥ 10–27
= 2.65 ¥ 10–7 kg
The net amount of deuterium consumed per day
= 2.65 ¥ 10–7 ¥ 24 ¥ 60 ¥ 60
= 0.02287 kg
= 0.023 kg
E xamplE 29 A cyclotron with Dee’s of diameter 1.8 m has a magnetic field of 0.8 Tesla. Calculate the energy
to which the doubly ionised helium ion He++ can be accelerated. Also calculate the number of revolutions the
particle makes in attaining this energy. Mass of He++ is 6.68 ¥ 10–27 kg.
Solution Given B = 0.8 T and mass of a-particle = 6.68 ¥ 10–27 kg, charge on a-particle (qa) = 2 ¥ 1.6 ¥ 10–19 C and r
= 0.9 m.
B 2 qa2 r 2
Maximum energy attained E =
2m
(0.8) 2 ¥ (3.2 ¥ 10-19 )2 ¥ (0.9)2
=
2 ¥ 6.68 ¥ 10 -27
= 0.39734 ¥ 10 -11 J
0.39734 ¥ 10 -11
= = 24.83 ¥ 106 eV
1.6 ¥ 10 -19
= 24.83 MeV
Frequency can be obtained by using the relation
Bqa 0.8 ¥ 3.2 ¥ 10 -19
f = =
2p m 2 ¥ 3.14 ¥ 6.68 ¥ 10 -27
= 0.061 ¥ 108
= 6.1 ¥ 106 Hz
Hence number of complete revolution performed by helium ion in obtaining the above energy
f 6.1
= = ¥ 106
2 2
= 3.05 ¥ 106 per sec
E xamplE 30 A cyclotron has an oscillator frequency of 12 ¥ 106 Hz and Dee radius of 21 inches. What is the
value of magnetic induction needed to accelerate deuteron in it?
Solution Given f = 12 ¥ 106 Hz, r = 21 inch = 0.53 m,
qd =1.6 ¥ 10–19 C and md = 2mp= 3.34 ¥ 10–27 kg
Formula used is 2p fmd 2 ¥ 3.14 ¥ 12 ¥ 106 ¥ 3.34 ¥ 10 -27
B= =
qd 1.6 ¥ 10 -19
= 1.573 T
E xamplE 31 Deuteron in a cyclotron describes a circle of radius 0.32 m just before emerging out of the
Dee’s. The frequency of the applied e.m.f. is 10 MHz. Find the flux density of the magnetic field and the
velocity of the deuterons emerging out of the cyclotron. Mass of deuteron is 3.32 ¥ 10–27 kg and charge 1.6
¥ 10–19 C.
Solution Given md = 3.32 ¥ 10–27 kg, qd = 1.6 ¥ 10–19 C,
f = 10 ¥ 106 Hz and r = 0.32 m.
Formula used is
2p md f 2 ¥ 3.14 ¥ 3.32 ¥ 10 -27 ¥ 10 ¥ 106
B= =
qd 1.6 ¥ 10 -19
= 1.303 T
Velocity of deuteron is
mv 2 qBr
= qvB or v =
r md
1.6 ¥ 10 -19 ¥ 1.303 ¥ 0.32
v= = 2.009 ¥ 107 m/sec
3.32 ¥ 10 -27
= 2.01 ¥ 107 m/sec
ExamplE 32 A betatron working on an operating frequency of 60 Hz has a stable orbit of 1.6 m diameter. Find
the energy gained per turn and also the final energy if the magnetic field at the orbit is 0.5 Tesla.
Solution Given r = d/2 = 0.8 m, f = 60 Hz and B = 0.5 T.
Average energy gained by the electron per turn is 4ewr2B. This can be proved as follows.
Let us consider that magnetic flux in the betatron is given by the relation
f = f0 sin wt
The increasing magnetic flux is obtained during the quarter cycle for a given direction in which the current in the
electromagnet increases from zero to maximum value.
T 2p p
\ Time of acceleration = = =
4 4w 2w
Where T is the time period of the changing magnetic flux and w is the corresponding angular frequency.
Energy gained by the electron per turn = eE
df d
=e = e (f0 sin w t )
dt dt
d
= ef0 (sin w t )
dt
T p
As this energy is gained in a time =
4 2w
Average value of energy per turn
p / 2w
ef0 d
=
p / 2w Ú dt ¢
(sin w t )dt ¢
0
2ewf0
=
p
To maintain a stable orbit of constant radius, tangential force on the electron must be zero. From this condition, we get
f0 = 2p r 2 B
2ew
\ Average energy per turn = ¥ 2p r 2 B
p
= 4ew r 2 B Joule
4 ¥ 1.6 ¥ 10 -19 ¥ 2 ¥ 3.14 ¥ 60 ¥ (0.8) 2 ¥ 0.5
= eV
1.6 ¥ 10 -19
= 482.3 eV
The total (or final) energy

= number of revolution done ¥ average energy per revolution
c
= ¥ 4ew r 2 B
4w r
= cerB
= 3 ¥ 108 ¥ 1.6 ¥ 10 -19 ¥ 0.8 ¥ 0.5 J
= 1.92 ¥ 10 -11 J
= 120 MeV
E xamplE 33 In a 70 MeV betatron synchrotron, the radius of the stable electron orbit is 28 cm. Find the
value of magnetic field B at the orbit for the given energy.
Solution Given E = 70 ¥ 106 eV = 70 ¥ 1.6 ¥ 10–13 J and r = 0.28 m.
E 70 ¥ 1.6 ¥ 10 -13
Formula used is B= =
cer 3 ¥ 108 ¥ 1.6 ¥ 10 -19 ¥ 0.28
= 0.83 T.
E xamplE 34 A sample of uranium emitting a-particles of energy 4.18 MeV is placed near an ionisation
chamber. Assuming that 12 particles per second enter the chamber, calculate the current produced, if an ion
pair requires energy of 40 eV. Charge on the electron e = 1.6 ¥ 10–19 C.
Solution Given Ea = 4.18 MeV = 4.18 ¥ 106 eV.
Energy required to produce an ion pair = 40 eV
Number of a-particles entering the chamber per second = 12
\ Energy required = 12 ¥ 4.18 ¥ 106 eV
= 50.16 ¥ 106 eV
Therefore, number of ion pair produced per second
Total energy supplied to the system
( n) =
Energy required to produce one ion pair
50.16 ¥ 106 eV
n=
40 eV
n = 1254 ¥ 103
Current (i) = time rate of collection of charge = 1254 ¥ 103 ¥ 1.6 ¥ 10–19 A
= 2.0 ¥ 10–13 A
E xamplE 35 A GM counter collects 108 electrons per discharge when the counting rate is 500 counts per
minutes. What will be the average current in the circuit?
500
Solution Number of discharge per second = 8.333
60
The average current in the circuit (i)
= 8.333 ¥ 108 ¥ 1.6 ¥ 10–19
= 1.33 ¥ 10–10 A
E xamplE 36 Neglecting parallel component of velocity, calculate cyclotron frequency and Larmour radius
for a 10 keV electron in the earth’s magnetic field of 5 ¥ 10–5 Tesla.
Solution Given B = 5 ¥ 10–5 T and E = 10 keV = 104 eV
Cyclotron frequency (w c ) = qB
m
1.6 ¥ 10 -19 ¥ 5 ¥ 10 -5
=
9.1 ¥ 10 -31
= 0.879 ¥ 107 sec -1
1/ 2
1 2 È 2E ˘
1/ 2
È 2 ¥ 104 ¥ 1.6 ¥ 10 -19 ˘
E= mv^ or v^ = Í ˙ =Í ˙
2 Îm˚ Î 9.1 ¥ 10 -31 ˚
= 5.93 ¥ 107 m/s
7
v^ 5.93 ¥ 10
and Larmour radius (rL) = =
w c 0.879 ¥ 107
= 6.746 m
E xamplE 37 What would be the cyclotron frequency of a solar wind proton streaming under the effect of
magnetic field B = 5 ¥ 10–9 Tesla. If the proton streams with velocity 3 ¥ 105 m/s, what would be the Larmour
radius. Neglect the parallel component of velocity.
Solution Given B = 5 ¥ 109 T and v^ = 3 ¥ 105 m/sec
mv^ 1.67 ¥ 10 -27 ¥ 3 ¥ 105
Larmour radius(rL) = =
eB 1.6 ¥ 10 -19 ¥ 5 ¥ 10 -9
= 6.26 ¥ 105 m
E xamplE 38 A He+ ion of energy 1 keV is gyrating in a circle of Larmour radius of 0.188 m under the effect
of external magnetic field. Calculate the magnetic field B by neglecting the parallel component of velocity.
Solution Given E = 103 eV and rL = 0.188 m.
1/ 2
1 2 È 2E ˘
Energy E = mv^ or v^ = Í ˙
2 Îm˚
1/ 2
È 2 ¥ 103 ¥ 1.6 ¥ 10 -19 ˘
v^ = Í -27 ˙
Î 4 ¥ 1.67 ¥ 10 ˚
= 2.19 ¥ 105 m/sec
and
mv^ 4 ¥ 1.67 ¥ 10 -27 ¥ 2.19 ¥ 105
B= =
erL 1.6 ¥ 10 -19 ¥ 0.183
= 4.996 ¥ 10 -2 T
E xamplE 39 Calculate the Larmour radius for a 3.5 MeV He++ ash particle in an 8 T DT fusion reactor by
neglecting the parallel component of velocity.
Solution Given E = 3.5 ¥ 106 eV.
1/ 2
1 2 È 2E ˘
Energy E = mv^ or v^ = Í ˙
2 Îm˚
1/ 2
È 2 ¥ 3.5 ¥ 106 ¥ 1.6 ¥ 10 -19 ˘
=Í ˙
Î 4 ¥ 1.67 ¥ 10 -27 ˚
V^ = 1.29 ¥ 107 m/s
and Larmour radius

mv^ 4 ¥ 1.67 ¥ 10 -27 ¥ 1.29 ¥ 107
rL = =
qB 2 ¥ 1.6 ¥ 10 -19 ¥ 8
= 3.36 ¥ 10 -2 m
E xamplE 40 Calculate the Debye length (lDe) and plasma frequency fPe for the plasma of earth’s ionosphere
having electron density of 1012 per m3 and thermal energy KTe as 0.1 eV.
Solution Given n = 1012 per m3 and KTe = 0.1 eV.
1/ 2
È e KT ˘
1/ 2
È 8.85 ¥ 10-12 ¥ 0.1 ¥ 1.6 ¥ 10-19 ˘
Debye length l De = Í 0 2 e ˙ = Í 12 -19 -19 ˙
Î ne ˚ Î 10 ¥ 1.6 ¥ 10 ¥ 1.6 ¥ 10 ˚
= 2.35 ¥ 10-3 m
1/ 2 1/ 2
È ne 2 ˘ È 1012 ¥ (1.6 ¥ 10-19 ) 2 ˘
w Pe = Í ˙ =Í -31 -12 ˙
Î me 0 ˚ Î 9.1 ¥ 10 ¥ 8.85 ¥ 10 ˚
w Pe = 0.5638 ¥ 108 rad/sec
w
So the plasma frequency ( f Pe ) = Pe
2p
0.5638 ¥ 108
or ( f Pe ) =
2 ¥ 3.14
= 0.08977707 ¥ 108 = 8.98 MHz
E xamplE 41 Calculate the plasma frequency and Debye length for a glow discharge of density 1016 per m3
and thermal energy 2 eV.
Solution Given density n = 1016 per m3 and KTe = 2 eV = 2 ¥ 1.6 ¥ 10–19 J
1/ 2
È e KT ˘
Debye length l De = Í 0 2 e ˙
Î ne ˚
1/ 2
Ê 8.85 ¥ 10 -12 ¥ 2 ¥ 1.6 ¥ 10 -19 ˆ
=Á ˜ = [11.0625 ¥ 10 -9 ]1/ 2
Ë 1016 ¥ (1.6 ¥ 10 -19 ) 2 ¯
= 1.0518 ¥ 10 -4 m
Angular plasma frequency
1/ 2
Ê ne 2 ˆ
(w Pe ) = Á ˜
Ë e 0m ¯
1/ 2
Ê 1016 ¥ (1.6 ¥ 10-19 ) 2 ˆ
=Á ˜
Ë 8.85 ¥ 10-12 ¥ 9.1 ¥ 10-31 ¯
= 5.637 ¥ 109 rad/sec
So plasma frequency
w Pe
f Pe =
2p
5.637 ¥ 109
=
2 ¥ 3.14
= 8.977 ¥ 108 Hz
Q.1 The radius of a nucleus depends on its mass number A and it is proportional to
(a) A (b) A1/3 (c) A2/3 (d) A3
Q.2 The magnitude of the spin angular momentum of a nucleon in the nucleus is
(a) h (b) h/2 (c) 0 (d) h2
Q.3 The total angular momentum of a nucleus having odd mass number A is
(a) 0 (b) h
(c) odd half-integral multiple of h (d) integral multiple of h
Q.4 The total angular momentum of a nucleus having even mass number A is
(a) 0 (b) h
(c) odd half-integral multiple of h (d) integral multiple of h
Q.5 Bohr magnetons mB is given by
eh eh eh eh
(a) mB = (b) mB = (c) mB = (d) mB =
2me 2m p 2me 2m p
Q.6 The magnetic moment of a proton is
(a) 0 (b) positive (c) negative (d) undefined
Q.7 Electric quadrupole moment of a nucleus is
(a) always zero (b) a measure of deviation of the nucleus from its
spherical symmetry
(c) simply the charge (d) undefined quantity
Q.8 Nuclear forces are
(a) short-range forces (b) long-range forces
(c) always repulsive (d) Coulomb forces between protons
Q.9 Yukawa gave a theory for nuclear forces based on the exchange of
(a) electrons between protons and neutrons (b) mesons between protons and neutrons
(c) electromagnetic photons (d) X-rays between protons and neutrons
Q.10 Binding energy curve is the curve between
(a) BE/A and A (b) BE and A
(c) BE/A and charge (d) BE/A and Coulomb energy
Q.11 Binding energy of light nuclei is small because
(a) more nucleons reside at the nucleus surface
(b) mass of the nucleus is small
(c) charge of the nucleus is small
(d) Coulomb force supercedes nuclear forces
Q.12 Nuclides lying on the stability curve are
(a) more stable (b) more unstable
(c) having large number of neutrons (d) having large number of protons
Q.13 As per nuclear shell model
(a) neutrons and protons move in the same orbit within the nucleus
(b) neutrons start moving in electronic shell lying near the nucleus
(c) protons are converted into neutrons and vice-versa
(d) neutrons and protons move in their separate orbits within the nucleus
Q.14 Nuclear magic numbers represent

(a) total mass number of the nucleus (b) difference in number of neutrons and protons
(c) the atomic magic numbers (d) even number of Z or N so that nuclear shells are
completely filled.
Q.15 Nuclei with nuclear magic numbers
(a) have always one proton and one neutron
(b) do not have any isotope
(c) have more number of isotopes than their neighbours
(d) have equal number of isotopes as their neighbours
Q.16 As per nuclear shell model, each shell
(a) is limited to a certain maximum number of neutrons or protons
(b) has equal number of neutrons and protons
(c) is half filled
(d) is two-third filled
Q.17 Nuclear shell model gives results which are in good agreement with experiments, when the potential
energy function V(r) is
(a) independent of r (b) rectangular
(c) oscillatory (d) a combination of square well potential and the
oscillator potential
Q.18 Nuclear drop model assumes
(a) drops of protons and neutrons
(b) the shape of nucleus as hollow sphere
(c) protons to remain on the surface of nucleus
(d) nucleus as a liquid drop containing neutrons and protons as its molecules
Q.19 In nuclear drop model
(a) drops of protons are converted into drops of neutrons
(b) surface energy is neglected
(c) some properties of nuclear forces are analogous to the properties of forces that hold a liquid drop
together
(d) protons and neutrons move freely
Q.20 A Curie is a standard unit of radioactivity. Its value is
(a) 109 disintegrations/sec (b) 2.7 ¥ 1010 disintegrations/s
10
(c) 3.7 ¥ 10 disintegrations/sec (d) none of these
Q.21 Radioactivity is the phenomenon associated with
(a) electron emission from atom (b) fission of nuclei
(c) transformation of nuclei (d) none of these
Q.22 The half-life period and radioactive decay constant are related as
1 ln 2
(a) T1/ 2 = (b) T1/ 2 =
l l
l
(c) T1/ 2 = 0.6931 (d) none of these
24
Q.23 The half-life of 11 Na is 15 hours. How long does it take for 93.75% of a sample of this isotope to
decay?
(a) 35 hr (b) 60 hr (c) 15 hr (d) 90 hr
Q.24 The activity of a certain radionuclide decreases to 15% of its original value in 10 days. Its half-life
would be
(a) 10 days (b) 5 days (c) 2.65 days (d) 3.65 days
Q.25 b-decay corresponds to
(a) an electron detached from atom’s outermost orbit
(b) the emission of proton from the nucleus
(c) electromagnetic wave pulse
(d) an electron emitted by the nucleus
Q.26 The end product of uranium series A = (Un + 2) is
(a) 206Pb (b) 208Pb (c) 207Pb (d) None of these
Q.27 The method(s) for determining the age of a sample is (are)
(a) cranium dating (b) carbon dating
Q.28 The radiation obtained from radioactive substance are
(a) a-rays (b) b-rays (c) g-rays (d) all of these
Q.29 Factors on which the range of a-particle depends are
(a) the initial energy of the particle (b) the ionisation potential of the gas
Q.30 Geiger Nuttal rule gives the range of
(a) a-particle (b) b-particle (c) b+-particle (d) g-rays
Q.31 Parity is not conserved in
(a) a-decay (b) g-decay (c) b-decay (d) none of these
Q.32 A long-lived excited nucleus is called
(a) isotone (b) isotope (c) isomer (d) none of these
Q.33 The time during which pulses are recorded but are of smaller duration in a GM Counter is called
(a) recovery time (b) resolving time (c) dead time (d) none of these
Q.34 The time during which pulses are not recorded in a GM Counter is called
(a) dead time (b) recovery time (c) resolving time (d) none of these
Q.35 Which scintillator is used for detection of a-particle in a scintillator counter?
(a) NaI (b) Zinc sulphide (c) Anthracene (d) None of these
Q.36 The radiation detector/detectors based on image formation is/are
(a) bubble chamber (b) Wilson’s cloud chamber
(c) nuclear emulsion charged (d) all of these
particle detector
Q.37 If s be the microscopic cross-section and n be the number of nuclei per unit volume, then the
microscopic cross-section is the product
(a) s n (b) s ndx (c) nd s dx (d) s nx
Q.38 The energy of the fast neutron is of the order of
(a) 1.0 MeV (b) above 1.2 MeV upto 10 MeV
(c) 1.0 eV (d) none of these
Q.39 Which particle cannot be accelerated by cyclotron?
(a) neutron (b) proton (c) deutron (d) a-particle
Q.40 In a nuclear fission reactor

(a) control rods are used to slow down the fast neutrons
(b) moderator is used to slow down the fast neutron
(c) coolant is used to slow down the fast neutrons
(d) none of these
Q.41 Plasma is
(a) state between third and fourth states of matter
(b) every ionised gas
(c) neutral gas with a few electrons
(d) ionised gas with ions and electrons in almost equal numbers
Q.42 In laboratory plasma, the electron temperature is usually
(a) 1 eV to 5 eV (b) 1 keV to 3 keV (c) 0 K (d) 1 MeV to 3 MeV
Q.43 If lDe and L be the Debye length and dimension of plasma, respectively, then the condition for
quasineutrality of plasma is
(a) lDe M L (b) lDe << L (c) lDe >> L (d) lDe = L/2
Q.44 For the purpose of particle acceleration, a linear accelerator uses
(a) a static electric field
(b) static magnetic field
(c) combination of electric and magnetic fields
(d) laser radiation
–
Q.45 A wave breaking field E in the plasma, whose frequency is wPe, is almost equal to
(a) cme wPe/e (b) ceme wPe (c) c2me wPe/e (d) 10 cme wPe/e
Q.46 In plasma wake field accelerator, the plasma wave is generated by
(a) electron beam (b) laser pulse
(c) static electric field (d) static magnetic field
Q.47 In a laser wake field accelerator a short laser pulse of pulse length L is used in a plasma of frequency
wpe/e and excites the wake wave of wavelength lp under the condition
(a) L = lp (b) L = lp/2 (c) L = lp/3 (d) c/L = wpe /2p
Q.48 In a self-phase modulation laser wake filed accelerator, the density of plasma should be such that the
laser pulse length L and wake length lp satisfy the relation
(a) L = lp (b) L = lp/2 (c) L > lp (d) L = lp/3
Q.49 In a laser beat wave accelerator, where two laser pulses of frequencies w1 and w2 are used in a plasma
of frequency wpe, the following condition should be satisfied
(a) w1+ w2 = wpe (b) w1 – w2 = wpe (c) w1+ wpe = w2 (d) w1+ 2w2 = wpe
Q.50 The binding energy of a nucleus is equivalent to the
(a) mass of protons (b) mass of nucleus
(c) mass defect of nucleus (d) none of these
Q.51 The fusion process is only possible at high temperature because
(a) nuclei disintegrate
(b) molecule disintegrate
(c) the nuclei get sufficient energy to overcome the Columbian attractive forces
(d) none of these
Q.52 Which of the following is not a particle accelerator?

(a) Betatron (b) Cyclotron
(c) Scintillation counter (d) none of these
Q.53 Which of the following nuclear reaction is known as exothermic when its Q value is
(a) negative (b) zero (c) positive (d) none of these
Q.54 Which of the following nuclear reaction is known as endothermic when its Q value is
(a) negative (b) positive (c) zero (d) none of these
Q.55 Most of the suitable material is used as moderator in a nuclear reactor
(a) cadmium (b) heavy water (c) slow neutrons (d) none of these
Q.1 Define radioactivity.

Q.2 Describe properties of a, b and g-rays.
Q.3 What is decay constant? How is it related to the decay probability per nucleus per second?
Q.4 What types of radiations are emitted in radioactive disintegration?
Q.5 Compare the properties of a, g and g-rays.
Q.6 What is mean life of radioactive isotope?
Q.7 What do you mean by half-life?
Q.8 What is half-life of radium, uranium and polonium?
Q.9 What do you mean by activity of a radioactive substance?
Q.10 What is basic mechanism of detection of radiation? Explain briefly.
Q.11 Which is the most sensitive and accurate method used for the deflection of nuclear radiation?
Q.12 Distinguish between ionisation chamber and GM counter.
Q.13 Briefly explain solid-state detector and radiation detectors.
Q.14 Define neutron cross-section. What are its units?
Q.15 What is fission?
Q.16 What do you understand by liquid-drop model?
Q.17 How can a fission chain reaction be controlled?
Q.18 What is nuclear fusion?
Q.19 What do you understand by ignition temperature?
Q.20 What is Lawson criterion?
Q.21 Differentiate between inertial and magnetic confinements?
Q.22 What are charged particle accelerators?
Q.23 Discuss plasma wake field accelerator in short.
Q.24 Discuss laser accelerator in short.
Q.25 Write short-note on
(a) Radioactivity (d) Neutron cross-section
(b) Nuclear fission and its applications (e) Cloud chamber
(c) Nuclear fusion (f) Bubble chamber
(g) Radiation detectors (k) Nuclear reactors

(h) GM counter (l) Reactor criticality
(i) Scintillation counter (m) Lawrence cyclotron
(j) Solid state detectors (n) Betatron
P ractice P roblems
general questions
Q.1 Discuss basic properties of a nucleus in detail.
Q.2 Write a note on angular momentum of a nucleus.
Q.3 Discuss magnetic and electric properties of a nucleus.
Q.4 What do you understand by parity of a nucleus.
Q.5 Discuss charge independence property of nuclear forces.
Q.6 Discuss the meson theory of nuclear forces.
Q.7 Write down the correlation between binding energy and stability of nuclei.
Q.8 Write down the facts of nuclear shell model.
Q.9 Discuss theory of nuclear shell model.
Q.10 What are the applications of nuclear shell model?
Q.11 Discuss nuclear magic numbers and their significance.
Q.12 Brief out the nuclear liquid drop model.
Q.13 What are various terms that contribute to the calculation of binding energy of nucleus?
Q.14 Discuss the volumes and surfaces energies used in the nuclear drop model.
Q.15 Write down the facts of nuclear liquid drop model and the semiempirical binding energy formula.
Q.16 Discuss in brief the way out for determining various constants appeared in nuclear liquid drop model.
Q.17 What is natural radioactivity? Explain what is radioactive disintegration. State the laws of radioactive
decay and deduce them from first principles using probability concepts.
Q.18 What is mean life of a radioactive isotope? Show that the mean life is the time for nuclei to decay to
1/e times their original number.
Q.19 Define radioactive constant and half-life period. Prove that the radioactive constant of a substance is
the reciprocal of the time after which the number of atoms of the substance falls to 1/e of its original
value.
Q.20 Define mean life of a radioactive nuclide. Derive a relation between mean life time and radioactive
constant.
Q.21 Define half-life and radioactive nuclide. Derive a relation between half-life and radioactive constant.
Q.22 What is the difference between half-life and mean life in radioactivity?
Q.23 What is the cause of radioactivity? Give various types of radioactive decays and discuss the process
involved in all these decays?
Q.24 What are a-particles? How will you show experimentally that a-particle is an ionised helium atom?
Q.25 State the conditions for a-decay and explain why in a-decay of a radioactive nuclide the kinetic
energy of the emitted a-particle is little less than the disintegration energy?
Q.26 What is b-decay? Is b-spectrum discrete or continuous?

Q.27 Show that the law of conservation energy and momentum are not obeyed in b-decay. Show how
neutrino hypothesis explains this discrepancy and accounts for continuous b-ray spectrum.
Q.28 Discuss gas ionisation based radiation detectors, suggest how one can detect thermal neutrons?
Q.29 What are radiation detectors based upon ionisation phenomenon. Explain the phenomenon and discuss
in detail any one of the ionisation based detectors. State briefly the applications also.
Q.30 Give three properties of nuclear radiation used in detection instruments. What are gas filled ionisation
based nuclear detectors?
Q.31 Discuss the principle and technical details of the working of ionisation method for detecting radiation.
Q.32 Describe the construction, principle and working of an ionisation chamber. Explain the difference
between ionisation chamber and GM counter. Why is an ionisation chamber less sensitive to b-particle?
Q.33 Describe the construction, working and principle of GM counter.
Q.34 Discuss in detail the construction and working of GM counter. What is the dead time and recovery
time? What do you mean by quenching of a GM counter? What is its necessity? How is it achieved
internally?
Q.35 Describe briefly operation and testing of a GM counter. What is plateau as applied to a GM counter?
Give its applications. What are its limitations?
Q.36 Discuss the working and principle of a scintillation counter.
Q.37 Explain the principle, construction and operation of a scintillation radiation counter. Draw the block
diagram to show the main components. Briefly describe a few types of scintillators used. In what way
a scintillation counter is superior to a GM counter?
Q.38 Describe construction and working of a scintillation counter.
Q.39 Describe the construction and working of a cloud chamber.
Q.40 Describe Wilson’s cloud chamber. How does a cloud chamber make visible the path of an incoming
particle? What is its main advantage.
Q.41 Discuss the working principle of a cloud chamber. How is it different from a bubble chamber? Explain.
Q.42 What is a bubble chamber? Describe the process by which ionising particles make a visible track in a
bubble chamber. Give its advantages and drawbacks.
Q.43 Discuss the construction and working of a semiconductor radiation counter. Give a few advantages of
semiconductor detector.
Q.44 Define and explain the term nuclear reaction cross-section. What are its units? If a beam of N0 particle
is incident on a slab of thickness x of material, how many particles will emerge out of unit slab. Given
that the slab contains n atoms per unit volume and s is the cross-section of the reaction.
Q.45 What do you understand by nuclear fission? Explain the release of energy during nuclear fission.
Q.46 Give the main assumptions of liquid drop model of the nucleus. Justify the name liquid-drop model.
Q.47 Obtain the expression for the binding energy of a nucleus based on liquid-drop model.
Q.48 How can the fission chain reaction be controlled?
Q.49 Define chain reaction. What is nuclear reactor? State its principle. Explain the function of moderator
in uranium fed nuclear reactors.
Q.50 What is nuclear fission? Explain why the fission of nucleus like 235U, 233U with slow neutrons is the
only fission reaction that is used for producing energy?
Q.51 Distinguish between fission and fusion. Describe the principle of construction and working of a
nuclear reactor.
Q.52 Describe the phenomenon of nuclear fission. Explain nuclear fission on the basis of liquid-drop model.
Q.53 Explaining the use of absorbers and methods of enrichment of 235U. Give the construction, working
and applications of nuclear reactor.
Q.54 Explain the term thermonuclear energy or nuclear fusion. Discuss its importance in universe. Where
do the sun and other stars get their energy from?
Q.55 Describe a nuclear reactor. How does it work?
Q.56 What are similarities and dissimilarities between nuclear fission and fusion?
Q.57 Explain the terms: neutron cross-section, reactor criticality and shielding.
Q.58 Explain carefully the principle of linear accelerator. Deduce the expression for the energy of the
particle and length of cylinders in terms of the constants of the apparatus.
Q.59 (a) What is the difference between linear and circular accelerator?
(b) Which accelerator makes use of electromagnetic radiations for accelerating particle?
Q.60 Describe the principle, construction and working of a cyclotron. Derive expression for the maximum
kinetic energy achieved by a particle of mass m in terms of the applied magnetic field and Dee radius.
Also state the relation in terms of the frequency of the applied electric field. Discuss its limitations.
Q.61 Can a cyclotron be used to accelerate electrons? If not why?
Q.62 What is a betatron? Derive the betatron condition for successful acceleration of electrons. Briefly
describe its principle, construction and function of alternating magnetic field in it.
Q.63 What do you understand by plasma? Explain its quasineutrality and collective behaviour.
Q.64 What is plasma frequency? How does it depend on plasma density? Is it same for both the constituents
of the plasma?
Q.65 What is Debye length? Why do you need Debye length to be much smaller than the dimension of the
plasma?
Q.66 What are plasma-based particle accelerators? Name any three of them.
Q.67 Discuss plasma wake field accelerator. How is it different from laser wake field accelerator?
Q.68 Explain plasma beat wave accelerator. What are its merits and demerits compared with laser wake
field accelerator?
Q.69 What do you understand by self-modulated laser wake field accelerator? Why do you need a dense
plasma for the successful operation of this accelerator?
Q.1 The linear absorption coefficient m of lead for 1 MeV gamma rays is 0.74 cm calculate (a) half-
thickness of lead for these g-rays, and (b) thickness of lead required to reduce the intensity of g-rays to
1
of its original value. [Ans: (a) 0.94 cm (b) 9.32 cm]
1000
Q.2 One mg of radioactive material with half-life of 1600 years is kept for 2000 years. Calculate the mass
which would have decayed by this time. [Ans: 0.50 mg]
Q.3 The half-life of a radioactive substance is 2.5 days. Calculate the percentage of original material left
after 7.5 days. [Ans: 12.5%]
Q.4 One gram radioactive radium-226 decays with a half-life of 1620 years. Calculate decay constant and
mean life? [Ans: 1.36 ¥ 10–11 per sec, 2337.3 yrs]
Q.5 Calculate the activity of 1 mg radium - 226 which has a half-life of 1620 years.
[Ans: 0.98 milli Curie]
Q.6 The half-life of 238U against a-decay is 4.5 ¥ 109 yrs. Find the activity of 1.0 kg of 238U.
[Ans: 0.334 milli Curie]
Q.7 When a nucleus of 7Li is bombarded with a proton, two a-particles are formed. Calculate the kinetic
energy of the a-particle assuming negligible energy of the bombarding proton. [Ans: 8.67 MeV]
235
Q.8 A reactor is producing energy at the rate of 1500 kW. How many atoms of U undergo fission per
second? How many kg of 235U would be used in 1000 hours of operation assuming that on an average
energy of 200 MeV is released per fission? [Ans: 65.86 ¥ 10–3 kg]
4
Q.9 A cyclotron has a magnetic field of 10 Gauss and a radius of 80 cm. Calculate the frequency of the
alternating electric field that must be applied and to what energy deutrons can be accelerated? Mass of
deuteron = 2 a.m.u. [Ans: 15.4 MeV]
Q.10 A cyclotron oscillator frequency 1 MHz is used to accelerate protons. If the radius of the Dees be 60 cm,
what would be the magnetic field in Tesla? [Ans: 6.56 T]
Q.11 A GM counter with dead time of 300 ms records 16000 counts per minute. What is the dead time loss
in counting rate. [Ans: 5 ¥ 10–6 min]
Crystal Structure 14
Learning Objectives
L0 1 Understand crystalline, amorphous LO 5 Discuss packing fraction for sc, bcc, fcc,
solids, primitive lattice and Wigner- diamond, hcp, interatomic attractive/
Seitz primitive cell, and types of repulsive forces
crystals LO 6 Explain ionic bond, covalent bond,
LO 2 Know about translation vectors, metallic bond, van der Waals bond,
lattice planes, and significance and hydrogen bond, crystal structure
representation of Miller indices analysis i.e., Bragg’s law and
LO 3 Illustrate structures of NaCl, CsCl, spectrometer, Laue method, powder
and diamond, coordination number method
of simple cubic lattice, bcc lattice, fcc LO 7 Evaluate vacancies, concentration of
lattice Schottky defects and Frenkel defects,
LO 4 Learn about interplanar spacing and compositional and electronic defect
nearest neighbour distance and atomic
radius
Introduction
A crystal structure is a unique arrangement of atoms. It consists of a set of atoms which are identical in
composition, arrangement and orientation, called basis and a lattice. Bases are located upon the points of
a lattice, which is an array of points repeating periodically in three dimensions. The points can be thought
of as forming identical tiny boxes, called unit cells, that fill the space of the lattice. The lengths of the
edges of a unit cell and the angles between them are called the lattice parameters. A crystal structure and
symmetry play an important role in determining many of its properties, like electronic band structure and
optical properties.
It is clear that a crystal structure is formed by the addition of a basis of atoms to every lattice point.
Mathematically, it can be represented as
Crystal structure = Lattice + Basis
14.1 Types Of sOLids LO1

The solids are broadly classified into two groups, namely, crystalline solids and amorphous solids.
14.1.1 Crystalline solids
Crystalline solids are arranged in fixed geometric patterns or lattices. Ice, methanol and sodium chloride are
a few examples of crystalline solids. They have orderly arranged units and are practically incompressible.
Crystalline solids also show a definite melting point and so they pass rather sharply from solid to liquid state.
There are various crystalline forms which are divided into seven crystal system or shapes. They are cubic,
tetragonal, hexagonal, orthorhombic, monoclinic, trigonal and triclinic. The units that constitute these sys-
tems can be atoms, molecules or ions. Ionic and atomic crystals are hard and breakable with high melting
points.
14.1.2 Amorphous solids
A rigid material whose structure lacks crystalline periodicity is called an amorphous solid. It means the pattern
of its constituent atoms or molecules does not repeat periodically in three dimensions. Even amorphous
materials have some short range order at the atomic length scale due to the nature of chemical bonding. They
are considered supercooled liquids in which the molecules are arranged in a random manner somewhat as in
the liquid state. Glass and plastic are the examples of amorphous solids. Unlike crystalline solids, amorphous
solids do not have definite melting points.
14.2 UniT CeLL LO1

The smallest portion of a space lattice which can generate the complete crystal by repeating its own dimensions
in various directions is called a unit cell. A unit cell is defined by the length of its edges and by the angles
between them, as shown in Fig. 14.1.
y
g
a
x
� b
c
z
Figure 14.1
  
Vectors a , b and c are called lattice vectors that form primitive axes in the crystal structure. We also call
them crystallographic axes, as the directions defined by these vectors are nothing but crystal axes. These vec-
tors are used in translation vector and hence are called fundamental translation vectors. The magnitudes of
  
vectors a , b and c are called lattice constants that specify the distances of the bases along the crystal axes.
14.2.1 primitive Lattice Cell

  
The parallelepiped defined by primitive axes a , b , c is called a primitive cell. A primitive cell is a type of
unit cell which fills all the space under the action of suitable crystal translation operation. A primitive cell is
a minimum volume unit cell, as shown in Fig. 14.2. There is a density of one lattice point per primitive cell.
 
The volume of a primitive cell is defined by primitive axes a, b , c as ( a ¥ b ) ◊ c .
b
c
a
14.2.2 Wigner–seitz primitive cell

The smallest volume enclosed by the normal lines drawn from midpoints of the lines which connect lattice
point to all nearby lattice points is called a Wigner–Seitz primitive cell. It is shown in Fig. 14.3.
14.3 Types Of CrysTALs LO1

All crystals are classified into seven crystal systems on the basis of the shape of the unit cells. Bravais in 1948
explained that there are fourteen different types of crystal lattices under the seven crystal systems. These
seven types of crystal systems are tabulated in Table 14.1, and shown in Fig. 14.4.
Table 14.1
Relative Between Number of Possible
S.No. Name of System Examples
Primitives Angles Lattices
1. Cubic a=b=c a = b = g = 90° 3(P, F, I) Nacl
2. Trigonal a=b=c a = b = g π 90° 1(R) CaSO4
3. Tetragonal a=bπc a = b = g = 90° 2(P, I) NiSO4
4. Hexagonal a=bπc a = b = 90°, g = 120° 1(P) Quartz
5. Orthorhombic aπbπc a = b = g = 90° 4(P, C, F, I) KNO3
6. Monoclinic aπbπc a = g = 90°, π b 2(P, C) FeSO4
7. Triclinic aπbπc a π b π g π 90° 1(P) CuSO4
P = Primitive, C = Base centered, I = Body Centered, F = Face centered, R = Rhombohedral

a
c A
b a b
a a g
a
Triclinic
Simple Cubic Body Centered Cubic Face Centered Cubic (Primitive)
c c c
c B
a C
a a b
a a a a
a b a
a a b
a
Hexagonal Trigonal Tetragonal Tetragonal Monoclinic
(Primitive) (Rhombohedral) (Primitive) (Body Centered) (Primitive)
c
D
b
b
a
Orthorhombic Orthorhombic Orthorhombic Orthorhombic Monoclinic

(Primitive) (Base Centered) (Body Centered) (Face Centered) (Base Centered)
Figure 14.4
14.3.1 Cubic System
In a cubic system, there are three types of lattices, namely simple cubic, body centered cubic and face centered
cubic. In addition to these structures, other structures are also depicted in Fig. 14.4.
(a) Simple Cubic: It contains lattice points at all eight corners of the unit cell. It is represented by sc.
(b) Body Centered Cubic: It contains one additional lattice point at the centre of the body of the unit cell
including at all eight corners. It is represented by bcc.
(c) Face Centered Cubic: It contains lattice points at the centre of each face as well as at all eight
corners. It is represented by fcc.
B
14.4 TrAnsLATiOn VeCTOrs LO2
We take any lattice point O as an origin in a plane lattice shown
in Fig. 14.5. Any other point in the two-dimensional lattice can

be obtained by repeatedly translating the vectors a and b . These
vectors are known as basis vectors. Based on these basis vectors, O A
we obtain the plane lattice by their repeated translation. The
position vector of any other lattice point, i.e., translation vector,
can be represented as
   Figure 14.5
T = n1a + n2b
where n1 and n2 are the integers

 which represent the number of lattice points along the directions OA and OB,

respectively, and a and b are the primitives. In case of three-dimensional crystal structures, the arrangement
of points can be represented as
   
T = n1a + n2b + n3c
  
where a , b and c are the primitives along X, Y and Z axes, respectively.
14.5 LATTiCe pLAnes LO2

A crystal lattice is made of a large number of parallel (a) (c)
equidistant planes (Fig. 14.6) known as lattice planes

and can be chosen in a number of ways, as shown in
Fig. 14.6a, b, c and d.
14.6 MiLLer indiCes LO2

(b) (d)
The integers which determine the orientation of a crystal
plane in relation to the three crystallographic axes are
called Miller indices. In order to find the Miller indices, Figure 14.6
the reciprocals of the intercepts of the plane on the
axes in terms of lattice constants are reduced to the smallest integers in ratio. Miller indices are also called crystal
indices.
y
Let us assume that a, b and c are the magnitudes of fundamental translation
vectors along the three axes, respectively. Again we consider that ABC B
represents the plane whose Miller indices are to be obtained. Let OA, OB and
OC are the intercepts made by this plane along the three axes (Fig. 14.7).
In order to find the Miller indices, we do the following.
O x
(a) We find the intercepts along the three axes. A
(b) We express these intercepts as multiple of lattice parameters.

(c) We divide these intercepts by lattice parameters. C
(d) We take the reciprocal of these. z
(e) We clear these fractions by taking the LCM of the denominators and Figure 14.7
multiplying it. This gives the required Miller indices.
This can be understood in a better way in the following steps corresponding to the above steps a to e.
Directions x y z
Step (a) OA OB OC
Step (b) pa qb rc
pa qb rc
Step (c) or p q r
a b c
1 1 1
Step (d)
p q r
Step (e) (h k l)
14.6.1 Significance of Miller Indices

Same Miller indices are assigned to all parallel planes, for example, planes with coordinates x, y, z; 2x, 2y,
2z; –x, –y, –z; etc., are represented by the same Miller indices. The corresponding index for a plane, which
is parallel to one coordinate axis, is zero. The corresponding index is negative for a plane that cuts an axis on
the negative side of the origin. The negative index is represented by putting a bar above the index.
14.6.2 Representation of Miller Indices

Let us consider a simple cubic system (a = b = g = 90° and a = b = c). For this system in Fig. 14.8a, b, c, d,
e, f and g, we show Miller planes corresponding to (1̄ 0 0), (1 0 0), (0 1 0), (0 0 1), (1 1 0), (1 0 1) and (1 1 1)
planes in a cubic crystal respectively.
y y
y
H E H E (0 1 0)
H
_
G (1 0 0) F (1 0 0) F G
G
B A B A B
x x x
C D C C
D
z
z z
(a) (b) (c)
y y
E H E H
F G F G
A B A B
x x
D C D C (1 1 0)
(0 0 1)
z z
(d) (e)
y y
E E
H H
(1 1 1)
F G (1 0 1) F G
A B A B
x x
D C D C
z (f) z (g)
Figure 14.8
14.7 siMpLe CrysTAL sTrUCTUres LO3

Below we discuss simple crystal structures, for example, the structures of sodium chloride, cesium chloride
and diamond.
14.7.1 Structure of naCl
The sodium chloride structure is shown in Fig. 14.9. It consists of two
face centered cubic sublattices, one of Na ion having its origin at the point
(0, 0, 0) and the other of the Cl ion having its origin midway along a edges
of the cube say at the point ÊÁ , , ˆ˜ . The space lattice is therefore truly
1 1 1
Ë 2 2 2¯
fcc, with a basic of one Na ion and one Cl ion separated by one half the
body diagonal of a unit cube. There are four Na+ – Cl– ion pairs in each unit
cube, with different ions in the positions. Figure 14.9
The coordinates of Na ions are defined as

Ê 1 1 ˆ Ê 1 1ˆ Ê 1 1ˆ
(0,0,0), Á , ,0˜ , Á ,0, ˜ , Á 0, , ˜
Ë 2 2 ¯ Ë 2 2¯ Ë 2 2¯
The coordinates of Cl ions CI�
Ê 1 1 1ˆ Ê 1ˆ Ê 1 ˆ Ê 1 ˆ Cs�
ÁË , , ˜¯ , ÁË 0,0, ˜¯ , ÁË 0, ,0˜¯ , ÁË , 0,0˜¯
2 2 2 2 2 2
(0,0,0)
14.7.2 Cesium Chloride (CsCl) structure
In a CsCl crystal, we have a body centered cubic lattice, in which
the Cs atom is placed at the origin (0, 0, 0), i.e., at the body
center point and Cl atoms are placed at the corners of the bcc
Figure 14.10
Ê 1 1 1ˆ
lattice whose coordinates are given by Á , , ˜ . Structure of
Ë ¯
cesium chloride is shown in Fig. 14.10. 2 2 2
14.7.3 diamond structure
The space lattice of diamond is fcc. In diamond structure we have
two fcc lattices placed at (0, 0, 0) and ÊÁ 1 , 1 , 1 ˆ˜ which superimpose
Ë 4 4 4¯
each other. In diamond structure we have two carbon atoms placed
at (0, 0, 0) and ÊÁ , , ˆ˜ . The diamond crystal structure is shown in

1 1 1
Ë 4 4 4¯
Fig. 14.11. Figure 14.11
14.8 COOrdinATiOn nUMber LO3

In a crystal, every atom is surrounded by the other atoms. The number of nearest neighbours to the given
atom in crystal lattice is known as coordination number.
14.8.1 Simple Cubic Lattice 5

2
This simple cubic lattice contains lattice points at all
eight corners of the unit cell. Any corner atom has four 3
neighbours. In the same plane as well as other two 1
neighbours. One of them is exactly above and the other
is exactly below in the vertical plane (Fig. 14.12). Hence, 4
the coordination number in a simple cubic lattice is
6 (= 4+2). 6
14.8.2 Body Centered Cubic Lattice Figure 14.12

It contains one additional lattice point at the center of the
body of the unit cell in addition to those at eight corners. In this lattice, every centered atom is surrounded by
eight equidistant neighbours (Fig. 14.13). Hence, coordination number in a body centered cubic lattice is 8.
y
2 3
O
4
1
O
x
8
7
6 5
14.8.3 Face Centered Cubic Lattice

It contains lattice points at the centre of each face as well as at all eight corners. Any corner atom has 4 face
centered atoms of the surrounding in its own plane and 4 face centered atoms above and below this plane
(Fig. 14.14). Hence the coordination number for this case is 12 (= 4 ¥ 3).
14.9 inTerpLAnAr spACing LO4

Let us consider a set of planes (h k l) in a unit cell. This plane PQR makes intercepts a/h, b/k, c/l on the three
Cartesian co-ordinate axes X, Y and Z respectively, as shown in Fig. 14.15. Consider O as the origin at a lattice
point. Let the plane PQR (h k l) be parallel to the plane passing through the origin. Draw a perpendicular ON
from O (origin) to the plane. Thus ON = d, is the distance between adjacent planes or interplanar spacing
between parallel planes. The normal ON makes angles a, b and g with the x, y and z axes, respectively. From
Fig. 14.15 we obtain OP = a/h, OQ = b/k and OR = c/l.
ON d dh
Now from DONP, = cos a = =
OP a/h a
ON d dk y
From DONQ, = cos b = =
OQ b/k b
ON d dl
From DONR, = cos g = =
OR c /l c Q
Then according to the law of direction cosines,
we get
cos2 a + cos2 b + cos2 g = 1 d
b/ k b N
È h2 k 2 l 2 ˘
d2 Í 2 + 2 + 2 ˙ =1 O P
Îa b c ˚ c/ l g a /h x
or 1
d= a
È h2 k 2 l 2 ˘ R
Í 2 + 2 + 2˙
Îa b c ˚
For cubic crystal a = b = c, we get
a z
d=
2 2 2
[h + k + l ] Figure 14.15
For tetragonal crystal a = b π c, we get

1
d=
h + k2 l2
2
+ 2
a2 c
14.10 neAresT neighbOUr disTAnCe And ATOMiC rAdiUs LO4

The distance between the centers of two neighbouring atoms is called Q
r
nearest neighbour distance. For a closely packed crystal, this distance
would be 2r for an atom of radius r. The distance r is called atomic
radius, which can be represented in terms of edge of cube a for certain
a
unit cell structures. For example, for simple cubic (sc) structure the
diameter of an atom would be equal to the cube edge. That is 2r = a 2r
or r =a/2. However, in case of a body centered cubic (bcc) structure, S
shown in Fig. 14.16, the atomic arrangement is such that a
r
a 2 a
PQ = r + 2r + r
Since PR = RS = a, PS = a 2 P a R
Therefore, PQ = PS 2 + SQ 2 Figure 14.16
or (4r ) 2 = 2a 2 + a 2
a 3
or r=
4
For the face centered cubic structure (fcc), shown in Fig 14.17, one can S
easily obtain r
a
r= .
2 2
2r
a
14.11 pACking frACTiOn LO5 r
Atomic packing fraction is defined as the ratio of volume of atoms

occupying the unit cell to the volume of the unit cell relating to that r r
substance. It is also known as relative packing density. It is denoted by f.
We can calculate f for different cell structures. P a R
14.11.1 Simple Cubic (sc) Structure Figure 14.17
In a simple cubic cell structure, the number of atoms in unit cell = 1.

The atomic radius is given by half of the lattice constant, i.e., r = a/2.
2
Volume of atom occupying the unit cell = 1 ¥ p r 3 = p ÊÁ ˆ˜
4 4 a
3 3 Ë 2¯
3
Volume of the unit cell = a
3
4 Ê aˆ
pÁ ˜
Ë ¯ p
Therefore, atomic packing fraction f = 3 32 = = 0.52
a 6
It also means that the packing fraction is 52%.
14.11.2 Body Centered Cubic (bcc) Structure

In a body centered cubic lattice, the number of atoms per unit cell = 2
3
The atomic radius, r = a
4 3
4 4 Ê 3 ˆ
Volume of atoms occupying the unit cell = 2 ¥ p r 3 = 2 ¥ p Á a˜
3 3 Ë 4 ¯
Volume of the unit cell = a3
Therefore, the atomic packing fraction
3
4Ê 3 ˆ
2¥ Á
Ë 4 ˜¯
a
3 3
f = 3
= p = 0.68
a 8
or it is 68%.
14.11.3 Face Centered Cubic (fcc) Structure

In a face centered cubic lattice, the number of atoms per unit cell = 4
2
4
3
4 4 Ê 2 ˆ
Volume of atoms occupying the unit cell = 4 ¥ p r 2 = 4 ¥ p Á a˜
3 3 Ë 4 ¯
3
4 Ê 2 ˆ
4¥ pÁ a˜
Therefore, atomic packing fraction f = 3 Ë 4 ¯ = p = 0.74
a3 3 2
So, the atomic packing fraction is 74%.
14.11.4 Diamond Structure
In the diamond structure, the number of atoms per unit cell = 8
3
8 3
4 3 4 Ê 3 ˆ
Volume of atoms occupying the unit cell = 8 ¥ p r = 8 ¥ p ÁË a˜
3 3 8 ¯
3
4 Ê 3 ˆ
8¥ pÁ a˜
Therefore, atomic packing fraction f = 3 Ë 8 ¯ = 3 p = 0.34
a3 16
or the packing fraction is 34%.
14.11.5 Hexagonal Closed Packed (hcp) Structure

In the hexagonal cell, the number of atoms per unit cell = 6
a
The atomic radius, r =
2 3
4 3 4 Ê aˆ
Volume of atoms occupying the unit cell = 6 ¥ p r = 6 ¥ p ÁË ˜¯
3 3 2
3 3 2 8
Volume of the hexagonal unit cell = a c where c = a
2 3
3 3 2 8
\ Volume = a a
2 3
= 3 2a 3
3
4 Ê aˆ
6¥ pÁ ˜
Therefore, atomic packing fraction f = 3 Ë 2 ¯ = p = 0.74
3 2a 3 3 2
or the packing fraction is 74%.
14.12 pOTenTiAL energy CUrVe And nATUre Of inTerATOMiC fOrCes LO5

The interatomic forces between atoms of the solids are electrostatic in nature which can either be attractive
or repulsive.
(a) Interatomic Attractive Forces
These forces act between the atoms to hold them together. These forces are long range and act between the
atoms at a finite distance. These forces become zero for infinite interatomic distances.
(b) Interatomic Repulsive Forces
A repulsive force also acts between the atoms to prevent the merging with each other and keep them
at a finite distance. These forces are short-range and act only when the atoms come so close with each
other that their electron cloud overlap with one another.
To understand interatomic forces, we consider a pair of atoms A and B, which are separated by a distance r at
any time. Then, the potential energy due to force of attraction is given by
a
U attr = - (i)
rn
and the potential energy due to force of repulsion is given by
b
U rep = + (ii)
rm
The total potential energy of the system is therefore given by
(iii)
U = Uattr + Urep
a b
U =- n
+ m
r r
where a, b, n and m are the characteristic constants of the molecule. From the above equation it is clear that
if r = •, then U = 0.
The force acting between these two atoms can be written as
dU d È a b ˘
F =- = - Í- n + m ˙
dr dr Î r r ˚
- na mb (iv)
or F= +
r n +1 r m +1
For the fixed value of a, b, n and m, it is clear from this equation that the force F = 0 at a particular distance
r = r0. The potential energy and interatomic force curves are shown in Fig. 14.18.
When the distance r between the two atoms is very large, then it is clear from the Figs 13.18(a) and (b) that
no force acts between the atoms and hence the total potential energy is zero. If the atoms approach each
other and come close together nearly equal to the atomic diameter then repulsive force also begins to act and
attractive force also gets very large. If the atoms come even more closer to each other then the repulsive force
increases faster and the net force between atoms becomes repulsive in nature.
Repulsive Potential
Energy Total Force
Force
Potential Energy
Total Potential Energy

Repulsive Force
r = r0
r r = r0 r
Attractive Force
Attractive Potential Energy
(a) Potential Energy Curve (b) Interatomic Force Curve

Figure 14.18
At a particular distance (r = r0), the attractive and repulsive forces attain equal values and therefore no net
force acts between the atoms. This position is known as equilibrium position in which the molecule becomes
most stable because of the minimum potential energy.
14.13 differenT Types Of bOnding fOrCes LO6

The constituent particles of crystals have different types of charge distribution which provides different types
of binding forces. The binding forces in most cases are electrostatic in nature but the distribution of electrons
in various atoms are qualitatively different in different crystals. These binding forces are of different types, for
example, ionic bond, covalent bond, metallic bond, molecular bond (or van der Waals bonds) and hydrogen bonds.
14.13.1 ionic bond
The ionic bond is formed due to transfer of one or more electrons from one type of atoms that lose electrons
readily to the other type that have affinity for electrons. Due to the transfer of electrons, these atoms become
positive and negative ions. The two types of atoms, which are involved in ionic bonding, are of different
types. The arrangement of ions formed in the ionic bonding is such that the Coulomb attraction between
ions of opposite charges is stronger than the Coulomb repulsion between ions of the same charges. Thus, the
ionic bond results from the electrostatic interaction of oppositely charged ions. In this situation, both the ions
attract each other and form ionic bond, which is clearly shown in Fig. 14.19. The examples of ionic crystals
are NaCl, CsCl, KBr, KOH, etc. The ionic bond is non-directional in nature.
Attraction
+ + + +
Na (Z = 11) CI (Z = 17) Na+ CI–

Figure 14.19
The ionic bonds are strong, hard and brittle. The cohesive energy of ionic crystals is very high. Thus the ionic
crystals have very high melting point and possess high latent heat of fusion. The ionic crystals are insulators
in general because of their very low conductivity at ordinary temperature. The conductivity of ionic crystals
increase with increase in temperature. Many ionic crystals are soluble in water (polar liquid) but not soluble
in nonpolar liquids like ether.
14.13.2 Covalent bond
In covalent crystals, the valency electrons do not get transferred from one atom to another as happens in an
ionic bond but are shared equally by both of the atoms. This sharing of valence electrons of the constituent
atoms forms a covalent bond. The sharing of electrons takes place in such a way that an electron with spin up
pairs with an electron with spin down, if that electron can occupy the states as per Pauli’s exclusion principle.
The sharing of one pair of electrons forms a singlet covalent bond and a double bond is obtained when two
pairs of electrons are shared. The covalent bond is also known as a homopolar or electron pair bond. The
conductivity of a covalent bond is low and increases with temperature. These are directional bonds and are
very hard because the bond is very strong, for example in diamond.
The simplest example of covalent bonding is the H2
molecule, as shown in Fig. 14.20. As the hydrogen + + + + +
atoms come close to each other, each of the two
electrons are attracted by both the nuclei. In case of H H H2
an oxygen molecule, two oxygen atoms share two Z = 1 Z = 1
pairs of electrons, thus forming a covalent double Figure 14.20
bond. A triple bond is formed by the sharing of three
pairs of electrons in a nitrogen molecule. Covalent bonds are also formed between atoms of different elements
like HCl, H2O, NH3, etc.
14.13.3 Metallic bond
Metallic crystals are commonly known as metals. In the atoms of metals, the electrons in the outermost orbits
are loosely bound as the ionisation energy is low in the case of metals. These electrons are free to move
around among all the atoms and are called free electrons or conduction electrons. This way the metals have
residual positive ions. The electrostatic attraction between these positive ions and negative electron gas is
responsible for holding the solid together. This type of bonding is called metallic bonding.
The metallic bond is electrostatic in nature, though partially, and do not exert directional influence. The
metallic bond is weaker in nature than the covalent bond because of the fewer electrons bonding the nuclei.
However, it can be stronger for those metals in which the number of valence electrons is greater. They are good
conductors. Because of the presence of free electrons, they have high thermal and electrical conductivities.
Most of the atoms in the first four groups of the periodic table like Li, Na, Cu, Ag, Zn, Fe, etc. are good
examples in which metallic bond exists.
14.13.4 Molecular bond or van der Waals bond

Many solids or crystals are composed
of neutral atoms or molecules
without any transferring or sharing
of the electrons. Hence, the ionic,
covalent and metallic bonds are not
Figure 14.21
possible between these atoms. These substances remain bound by much weaker short-range attractive forces,
which are called van der Waals forces. These forces are weaker than the atomic bonding forces. The van der
Waals bonds are usually found in inert gases in which outermost electron orbits are completely filled, i.e.,
there are no valence electrons and hence they are incapable of forming any bond. These bonds are formed due
to electrostatic attraction between oscillating or permanent dipoles (Fig. 14.21). As we know that the dipoles
are formed due to the asymmetrical charge distribution around atoms, these dipoles are called oscillating
dipoles.
In the case of inert gases, there is a very small attraction between the atoms due to closed outer shells. These
gases condense when the temperature is reduced and hence a weak interatomic attraction is developed due
to van der Waals forces. The van der Waals forces are non-directional in nature and a little energy is required
to break the bonds because these are much weaker than ionic and covalent bonds. These types of bonds are
found in inert gases like solid argon and in many organic symmetrical molecules like methane (CH4).
14.13.5 hydrogen bond
In certain crystals, a positive hydrogen ion
O– –
(H+) attracts negative ions such as F–,O–,
N–, etc. Sometimes, due to electrostatic H+ H+
attraction, attachment between atoms in
different molecules or within a molecule O– – O– –
occurs in addition to bonds, which
holds atoms together to form molecules.
A hydrogen bond is formed when a H+ H+
H+ H+
hydrogen atom makes such an attachment
or association with an electronegative Figure 14.22
atom like oxygen, nitrogen, fluorine, etc.
The hydrogen bond is found in H2O, HF and in many organic molecules, particularly proteins and DNA
molecules. In the water molecule (H2O), the hydrogen and oxygen atoms are held together by covalent bonds.
The positive dipole end i.e., hydrogen, can strongly attract the negative dipole end of water molecule. This
bonding of the water molecule is shown in Fig. 14.22, where dashed lines represent the hydrogen bond.
14.14 CrysTAL sTrUCTUre AnALysis LO6

As we know that the X-rays can penetrate solids. Since the wavelength of the X-rays (~1 Å) is of the order of
interplanar spacing, these rays get strongly diffracted from different crystal planes. By analysing the different
X-rays, we can gather information about the internal atomic arrangement of a crystal.
14.14.1 Bragg’s Law: Diffraction of X-rays A C

The diffraction of X-rays from a single crystal is shown in B D
Fig. 14.23. When a narrow beam of X-rays of wavelength l
is incident on the parallel planes at an angle q, the X-rays are O
q q
scattered by the atoms of the crystals. Consider a set of parallel P
qq d
lattice planes which are separated by a distance d. Further, two N N¢
S
X-ray beams AO and BO¢ are incident on the crystal, as shown in O¢
the figure. The incident X-rays get diffracted at an angle q at the
Figure 14.23
point O and O¢ of the crystal planes. The path difference between these two beams can be obtained by
drawing perpendiculars ON on BO¢ and ON¢ on O¢D, as
NO¢ + O¢N¢ = d sin q + d sin q = 2d sin q
For constructive interference (or for maximum intensity), we must have
2d sin q = nl (i)
where n is an integer. The above relation is known as Bragg’s law of diffraction. It is useful in calculating
the distance d between crystal lattice planes once we know the wavelength l of the X-rays and measure the
angle of diffraction q.
It is clear from Bragg’s condition (i) that every X-ray will not get diffracted by the atoms of a crystal, only
those will be diffracted whose wavelength l and the angle q exactly match this condition. The standard
methods of X-ray diffraction used in the analysis of crystal structure are designed to achieve this. Bragg’s
X-ray spectrometer, Laue method and Powder method are such methods, which are discussed below.
14.14.2 Bragg’s X-ray Spectrometer

Bragg devised an apparatus used to study
the glancing angle q and the intensities of
diffracted X-rays for a given sample. The
essential parts of Bragg’s spectrometer are
shown in Fig. 14.24.
A monochromatic X-ray beam from the
source is passed through two slits, S1 and S2
which collimate it into a fine narrow beam.
This X-ray beam is then allowed to fall on the
crystal sample C mounted at the centre of a
turn table. A movable arm is also attached to
the turn table for detecting the reflected X-ray
beam from the crystal C. This turn table is
capable of rotation about a vertical axis and
Figure 14.24
the angle of rotation can be measured on the
circular scale. The rates of rotation of the
turn table and the detector (ionisation chamber) arm are such that the ionisation chamber always receives the
diffracted beam. When the turn table rotates through an angle q, the ionisation chamber arm automatically
rotates through an angle 2q with the direction of incident ray. This way the measurements of different diffracted
X-ray beams’ intensities and angles are recorded. Then using Bragg’s condition, we get the interplanar spacing
and the structure of the crystal.
14.14.3 Laue Method
The Laue method is useful for the determination of crystal structure. In this method, a single crystal is held
stationary in the path of an incident X-ray beam, as shown in Fig. 14.25.
When a continuous X-ray beam through a pinhole is allowed to fall on the crystal, then this beam is diffracted
by the crystal, the transmitted-diffracted and the reflected-diffracted beams are received by the films P and
Q, respectively as shown in (Fig. 14.25a).
Film Q Film P
Crystal
X-ray +
Beam
Pinhole
Collimator
Crystal Holder Laue Pattern
(a) (b)
Figure 14.25
The transmitted-diffracted beams form a series of spots, which is the characteristic of crystal structure and is
called Laue pattern (Fig. 14.25b). Each spot in the Laue pattern corresponds to the interference maxima for a
set of crystal planes satisfying the Bragg’s condition (2d sin q = nl) for a particular wavelength selected from
the beam of incident light. By studying the position and intensities of these Laue spots, the crystal structure
can be determined.
14.14.4 powder Method
It is a standard and straight forward
technique for analysing the crystal
structure. In this technique, we use the
crystal in powder form instead of single
crystal so that its tiny crystals (i.e.,
crystallities) are randomly (i.e., almost
continuously) oriented and make all
possible angles with the incident beam.
A small specimen of the crystalline
powder is taken in a small capillary
Figure 14.26
tube (P) of nondiffracting material and
is placed in the path of fine monochromatic beam of X-ray (Fig 14.26). Thus all possible diffraction planes
will be available for the Bragg diffraction (2d sin q = nl) to take place.
All these diffracted rays will lie on a conical surface hav- 2l
ing its apex at P and semivertical angle 2q. The diffracted
X-ray is recorded by the photographic film placed around
the crystal and we get the arc of the circle on the photo- Photo film
graphic film, as shown in Fig. 14.27. Figure 14.27
If 2l is the distance of the two arcs of the same circle on the

photographic film and R is the radius of cylindrical camera (Fig. 14.28), then
l l
2q = or q = C
R 2R 2q
R
2l
This value of q will given the spacing between the planes with the
help of Bragg’s relation,
Photo film
2d sin q = nl
Figure 14.28
By differentiating
2Dd sin q + 2d cos q ◊ Dq = 0
Dq tan q
or =-
Dd d
If q tends to 90° then 2q = 180°. For this angle, X-rays get reflected back along their initial path and such
reflected beams cannot be recorded. For q = 90°, (Dq/Dd) becomes very large so that small variations in d
produces large variation in q.
14.15 pOinT defeCTs in sOLids LO7

Point defect is a discontinuity in a crystal lattice. It consists of either a missing atom or an ion that creates a
vacancy in the lattice (often known as Schottky defect). If an extra atom or ion exists between two normal
lattice points it is said to create an interstitial position and if the missing atom or ion shifts to an interstitial
position, then the vacancy is called Frenkel defect. Point defect occurs because of the absence of a matrix
atom or the presence of an impurity atom or the matrix atom at the wrong place. The most common type point
defects in a pure crystal are given below.
14.15.1 Vacancies
Vacancies are created during crystallisation or from thermal
vibrations of the atoms at high temperatures. During thermal Vacancy
vibration, the atoms may acquire sufficiently high energy and
evaporate partially or completely and hence create a vacancy
(vacancies) in the lattice (Fig. 14.29).
In an ionic crystal, the formation of vacancy requires the charge
neutrality which should be maintained in the crystal as a whole. As a
result, a pair of vacancies causes to missing of one cation and one anion
from the structure. Such a pair of vacant sites is called a Schottky defect
(Fig. 14.30). If a cation goes into an interstitial position, then the
Figure 14.29
interstitially pair is known as Frenkel Defect (Fig. 14.30).
Evaluation of Concentration of Schottky Defects

In the process of Schottky defect, the + +
– – + – + – +
absorbed energy is more than compensated
by the resultant disorder in the lattice. Due – + – + – + – –
to this excess energy the disorder increases Schottky Defect
+ – + + – – +
that causes an enhancement in entropy. A
Frenkel
certain number of vacancies, called the Defect – + + – + –
concentration of Schottky defects, are –
+ – + – + – +
always present in thermal equilibrium
which can be calculated as follows. – + – + – + – + –
Let us consider that n is the number of + – + – – –
+ + +
Schottky defects produced in removing na+
cations and na– anions. Further, N is taken as Figure 14.30
the total number of cation-anion pairs. If U be the average energy, which is required to produce a Schottky defect,
then nU would be the increase in energy associated with the generation of n vacancies. The number of different
ways in which a cation or an anion can be removed is given by N! . Since there are n cation and n anions,
( N - n )!n!
the total number of different ways in which n Schottky defects can be produced will be
È N! ˘È N! ˘
W =Í ˙ Í ˙
Î ( N - n )!n!˚ Î ( N - n )!n !˚
The increase in entropy is given by
S = k ln W
È N! ˘
= 2k ln Í ˙
Î ( N - n )! n ! ˚
Here k is the Boltzmann constant. This increase in entropy produces a change in the Helmholtz free energy
F, which can be obtained as
F = Increase in energy – Temperature × Increase in entropy
È N! ˘
= nU - 2kT ln Í ˙
Î ( N - n )! n ! ˚
Using the Sterling’s approximation ln y! = y ln y – y, we get the following expression for the change in
Helmholtz free energy F
F = nU – 2kT[N ln N – (N – n) ln(N – n) – n ln n]
The energy should be a constant at equilibrium.
∂F È ( N - n) ˘
This will give = 0 or U = 2kT ln Í
∂n T Î n ˙˚
N -n Ê U ˆ
or = exp Á
n Ë 2kT ˜¯
For smaller number of Schottky defects, i.e., when n << N, we have N – n ª N.
Ê U ˆ
This gives, n = N exp Á -
Ë 2kT ˜¯
The above expression gives the number of Schottky defects in binary ionic crystals like MgO and NaCl
at ordinary temperature. It is clear from the expression that the fraction n/N of Schottky defects increases
exponentially with the temperature.
Evaluation of Concentration of Frenkel Defects

In order to calculate the concentration of Frenkel defects, we consider a pure crystal that consists of positively
and negatively charged ions in equal proportion. Further, we consider that it contains total N number of ions
and n Frenkel defects. The number of different ways in which n Frenkel defects will be produced is given as
N! Ni !
W= ◊ (i)
( N - n )!n! ( N i - n )!n !
where Ni is the number of interstitial sites. The Helmholtz free energy F of the crystal is given by the relation
F = U – ST (ii)
where S is the increase in entropy and T is the temperature.
If Ei be the energy required to produced a vacancy, then U can be expressed as
(iii)
U = nEi
and the associated increase in entropy is given by the Boltzmann relation
È N! Ni ! ˘
S = k ln W = k ln Í ◊ ˙ (iv)
Î ( N - n )!n! ( N i - n )!n !˚
With the help of Eqs. (iii) and (iv), Eq. (ii) becomes
È N! Ni ˘
F = nEi - kT ln Í ˙ (v)
Î ( N - n )!n! ( N i - n )!n !˚
Sterling’s approximation [ln y! = y! ln y – y] yields
È N! Ni ! ˘
ln Í
( N - n )! n ! ( N - n )! n !˙ = N ln N + Ni ln Ni - {( N - n ) ln( N - n)
Î i ˚
+ ( Ni - n ) ln( Ni - n ) + 2n ln n )
= nEi – kT[N ln N + Ni ln Ni – (N – n)ln(N – n)
– (Ni – n)ln(Ni – n) – 2n ln n] (vi)
Differentiating Eq. (vi) w.r.t. n, we get
È ∂F ˘ ( N - n )( Ni - n ) (vii)
ÍÎ ∂n ˙˚ = Ei - kT ln n2
T
Free energy remains constant, when the equilibrium position is attained at a given temperature T. It means
È ∂F ˘
ÍÎ ∂n ˙˚ = 0
T
( N - n )( Ni - n )
\ Ei - kT ln =0
n2
( N - n )( Ni - n ) Ei
or ln =
n2 kT
For smaller number of Frenkel defects, i.e., when N >> n and Ni >> n the above relation reads
NNi Ei
ln =
n2 kT
Ei
or 2ln n = ln( NNi ) -
kT
1 E
or ln n = ln( NNi ) - i
2 2kT
or n = ( NNi )1/2 e - Ei / 2 kT
From the above relation the concentration of Frenkel defects at a temperature T can be calculated.
14.15.2 interstitial
As defined earlier, in interstitial defect an atom or an ion moves from its proper position to a place between
regular lattice sites, as shown in Fig. 14.31. The interstitial is either due to a normal atom of the crystal or of
a foreign atom.
14.15.3 Compositional Defect
The compositional defect occurs because of the replacement of a host atom by a foreign atom. The foreign
atom remains at the regular lattice site, as shown in Fig 14.32.
14.15.4 Electronic Defect
At absolute zero in a purely covalent crystal (e.g. Si), Si Si Si
the electrons are tightly bound to the core and all are
Electronic Free Electron
said to be in the valence band. Above absolute zero, Defect
some of the electrons are likely to occupy higher
Si Si Si
energy state depending upon the temperature. So
in the crystal of pure silicon, some of the electrons
from the covalent bonds get thermally released Hole
and become free to move, as shown in Fig. 14.33. Si Si Si
This way the deficiency of electron creates a hole.
Then the electrons and holes give rise to electronic
imperfections. Figure 14.33
s UmmarY

✦ A unique arrangement of atoms is called a crystal structure. It consists of a set of atoms which are
identical in composition, arrangement and orientation called basis and a lattice. Bases are located upon
the points of a lattice, which is an array of points repeating periodically in three dimensions. A crystal
structure and symmetry play an important role in determining many of its properties such as electronic
band structure and optical properties.
✦ The solids are broadly classified into two groups, namely crystalline solids and amorphous solids.
Crystalline solids are arranged in fixed geometric patterns or lattices. Ice, methanol and sodium chloride
are a few examples of crystalline solids. A rigid material whose structure lacks crystalline periodicity
is called an amorphous solid. It means the pattern of its constituent atoms or molecules does not repeat
periodically in three dimensions in amorphous solids.
✦ A crystal structure can be obtained by translation of a unit cell in three dimensions. A unit cell is
a smallest pattern of a space lattice, which can generate the complete crystal by repeating its own
dimensions in various directions in amorphous solids.
✦ Based on the shape of the unit cells all crystals are classified into seven crystal systems. These are
cubic system, trigonal system, tetragonal system, hexagonal system, orthorhombic system, monoclinic
system and triclinic system. In a cubic system, there are three types of lattices called simple cubic (sc)
system, body centered cubic (bcc) system and face centered cubic (fcc) system.
✦ With reference to a lattice point as an origin in a plane lattice, any other point can be obtained by
 
repeatedly translating the vectors a , b and c , which are the primitives along X, Y and Z axes, respectively.
The position vector of this lattice point, i.e., translation vector, can be represented as
   
T = n1a + n2b + n3c
where n1, n2, and n3 are the integers which represent the number of lattice points along the three
directions.
✦ A crystal lattice is made of large number of parallel equidistant planes known as lattice planes.
✦ The integers which determine the orientation of a crystal plane in relation to the three crystallographic
axes are called Miller indices. In order to find the Miller indices, the reciprocals of the intercepts of
the plane on the axes in terms of lattice constants are reduced to the smallest integers in ratio. Miller
indices are also called crystal indices.
✦ In a crystal, every atom is surrounded by the other atoms. The number of nearest neighbours to the
given atom in the crystal lattice is known as coordination number.
✦ For a set of planes (h k l) in a unit cell, the distance between adjacent planes or interplanar spacing
between parallel planes is given by d = 1 . Here a, b and c are the fundamental
Èh 2
k2 l2 ˘
Í 2 + 2 + 2˙
Îa b c ˚
translational vectors along the three axes.
✦ The distance between the centers of two neighbouring atoms is called the nearest neighbour distance.
For a closely packed crystal, this distance is 2r for an atom of radius r. The distance r is called atomic
radius, which is generally represented in terms of edge of cube a for certain unit cell structures.
✦ The ratio of volume of atoms occupying the unit cell to the volume of the unit cell relating to that
substance is called atomic packing fraction. It is also known as relative packing density. It is denoted
as f. The atomic packing fraction f for the simple cubic structure is 52%. For a body centered cubic
(bcc) structure f = 68% and for a face centered cubic (fcc) and hexagonal closed packed structures f =
74%. The atomic packing fraction f for the diamond structure is 34%.
✦ The interatomic forces between atoms of the solids are electrostatic in nature which can either be
attractive or repulsive. It is obtained that at a particular distance, the attractive and repulsive forces
attain equal values and therefore no force acts between the atoms.
✦ The constituent particles of crystals have different types of charge distribution which provides different
types of binding forces. The binding forces in most cases are electrostatic in nature but the distribution
of electrons in various atoms are qualitatively different in different crystals. These binding forces are
of different types, for example, ionic bond, covalent bond, metallic bond, molecular bond (or van
der Waals bonds) and hydrogen bonds. Accordingly, the crystals are referred to as the ionic crystal,
covalent crystal, metallic crystal, molecular crystal and hydrogen bonded crystal.
✦ Since the X-rays can penetrate solids and their wavelength (1 Å) is of the order of interplanar spacing,
these rays can get strongly diffracted from different crystal planes. An analysis of the diffracted X-rays
can provide the information about the structure of the crystal. The standard methods of X-ray diffraction
include Bragg’s X-ray spectrometer, Laue method, rotating crystal method and powder method.
✦ Point defect is a discontinuity in a crystal lattice. It consists of either a missing atom or an ion that
creates a vacancy in the lattice (often known as Schottky defect). If an extra atom or ion exists between
two normal lattice points, it is said to create an interstitial position and if the missing atom or ion shifts
to an interstitial position, then the vacancy is called Frenkel defect. Point defect occurs because of the
absence of a matrix atom or the presence of an impurity atom at the matrix atom in the wrong place.
✦ The number of Schottky defects in binary ionic crystals like MgO and NaCl at ordinary temperature
is given by n = N exp ÊÁ -
U ˆ
where N is the total number of cation-anion pairs and U is the average
Ë 2kT ˜¯
energy required to produce the Schottky defects.
Ê E ˆ
✦ The number of Frenkel defects in crystals at ordinary temperature is given by n = ( NNi )1/ 2 exp Á - i ˜
Ë 2kT ¯
where N is the number of ions, Ni is the number of interstitial sites and Ei is the energy required to
produce the vacancy.
s olved e XamPles
E xamplE 1 A plane cuts intercepts 2a, 3b and c along the crystallographic axes in a crystal. Determine the
Miller indices of plane.
Solution Intercepts are 2a, 3b and c.
Then from the law of rational indices, we have
a b c
2a : 3b : c = : :
h k l
1 1 1
or : : = 2 : 3:1
h k l
1 1
or h:k :l = : :1 = 3: 2 : 6
2 3
Therefore, the Miller indices of the plane are (3 2 6).
E xamplE 2 In a triclinic crystal, a lattice plane makes intercepts at a length a, 2b and –3c/2. Find the Miller
indices of the plane.
Solution Intercepts are a, 2b and –3c/2.
-3c a b c
\ a : 2b := : :
2 h k l
1 1 1 -3
or : : = 1: 2 :
h k l 2
1 -2
or h : k : l = 1: : = 6 : 3: - 4
2 3
–
Therefore, the Miller indices of the given plane are (6 3 4).
  
E xamplE 3 Deduce the Miller indices for planes in each of the following sets which intercept a , b and c at
(i) 3a, 3b, 2c (ii) a, 2b, • (iii) a, b/2, c
Solution (i) Intercepts are 3a, 3b, 2c.
Then,
a b c
3a : 3b : 2c = : :
h k l
1 1 1
or 3:3:2 = : :
h k l
1 1 1
or h : k :l = : :
3 3 2
or h : k : l = 2:2:3
Therefore, the Miller indices are (2 2 3).
(ii) Intercepts are a, 2b, •.
Then,
a b c
a : 2b : • = : :
h k l
1 1 1
or : : = 1: 2 : •
h k l
1 1
or h : k : l = 1: : = 2 :1: 0
2 •
Therefore, the Millers indices are (2 1 0).
(iii) Intercepts are a, b/2, c.
Then,
b a b c
a: :c = : :
2 h k l
1 1 1 1
or 1: :1 = : :
2 h k l
or h:k:l=1:2:1
ExamplE 4 Calculate the spacing between (1 0 0) and (1 1 1) planes of a cubic system of lattice parameter a.
Solution Spacing between the planes of a cubic system of lattice parameter a.
a
d hkl =
h + k2 + l2
2
a
For plane (1 0 0), d100 = =a
1 + 02 + 02
2
a a
and For plane (1 1 1), d111 =
1 2
1 2
1 2 3
E xamplE 5 Deduce the Miller indices of a set of parallel planes which make intercepts in the ratio of a: 2b on
  
the x and y axis and are parallel to z-axis, a , b , c being primitive vectors of lattice. Also calculate the interplanar
distance d of the plane taking the lattice to be cubic with a = b = c = 5Å.
Solution The parallel planes are parallel to z axis. It means that their intercepts on the z-axis are infinite. Thus, the
intercepts are a, 2b and •. And, lattice constant a = b = c = 5Å.
Then,
a b c
a : 2b : • c = : :
h k l
1 1 1
or : : = 1: 2 : •
h k l
1 1
or h : k : l = 1: : = 2 :1: 0
2 •
\ Interplanar distance
a 5 ¥ 10-10
d= =
h2 + k 2 + l 2 22 + 12 + 02
5 ¥ 10-10
= = 5Å
5
E xamplE 6 Determine the Miller indices of plane parallel to the z axis and cut intercepts of 2 and 2/3 along
x and y axes, respectively.
Solution Intercepts are 2a, 2b , •.
3
2b a b c
2a : :• c = : :
3 h k l
1 1 1 2
or : : = 2: :•
h k l 3
1 3 1
or h:k :l = : : = 1: 3: 0
2 2 •
E xamplE 7 Calculate the interplanar spacing for (2 3 1) plane of an fcc structure whose atomic radius is
0.175 nm.
Solution Given plane = (2 3 1) and atomic radius (r) = 0.175 nm.
Atomic radius (r) of fcc structure
2
= a
4
a
and interplanar spacing d=
h2 + k 2 + l 2
2a 4r 4 ¥ 0.175 ¥ 10-9
 r= fia= =
4 2 2
a 4 ¥ 0.175 ¥ 10-9
\ d 231 = =
h2 + k 2 + l 2 2 ¥ 22 + 32 + 1
4 ¥ 0.175 ¥ 10-9
= = 0.132 10 9 m
2 ¥ 14
E xamplE 8 In a simple cubic crystal (i) find the ratio of intercepts of three axes by (1 2 3) plane and (ii) find
the ratio of spacings of (1 1 0) and (1 1 1) planes.
Solution (i) Given (h k l) of the plane as (1 2 3). Intercepts on the axes of a simple cubic crystal are given as a/h, a/k, a/l.
\ The ratio of intercepts are
a a a 1 1
: : = 1: :
1 2 3 2 3
(ii) The spacings d of plane (h, k, l) in a simple cubic crystal of side a
a
=
h + k2 + l2
2
a a
\ For plane (1 1 0), d110 = =
2
1 +1 + 0 2 2
a a
and for plane (1 1 1), d111 = =
2
1 +1 +1 2 2 3
Therefore the ratio of spacing between these two plane is d110 / d111 = 3/ 2 = 1.225
E xamplE 9 Calculate the distance between two atoms of basis of the diamond structure if the lattice constant
of the structure is 5Å.
Solution Given lattice constant a = 5Å.
The distance between two atoms is equivalent to the nearest neighbor distance.
3
For diamond structure, nearest neighbour distance = a
4
3 1.732 ¥ 5
\ Distance between two atoms = ¥ 5Å = = 2.17Å
4 4
E xamplE 10 What is the number of atoms in the primitive cell of diamond. Calculate the length of a primitive
translation vector if the cube edge a = 3.56 Å.
Solution Diamond is a fcc lattice with two carbon atoms in a primitive cell. So, the number of atoms is 8.
Given the cube edge a = 3.56 Å.
a 3.56
\ Primitive translation vector = = = 2.52Å
2 1.41
ExamplE 11 Determine the number of atoms per unit cell of lead which has an fcc structure. Atomic weight of
Pb = 207.2, density of Pb = 11.36 ¥ 103 kg m–3, and a = 3.2 Å and Avogadro’s number = 6.023 ¥ 1026/kg mole.
Solution Given atomic weight of Pb (M) = 207.2, density of Pb = 11.36 ¥ 103 kg m–3 and a = 3.2 Å
Avogadro’s number N = 6.023 ¥ 1026 /kg mole.
a3r N
Number of atoms n =
M
(3.2 ¥ 10-10)3 ¥ 11.36 ¥ 103 ¥ 6.023 ¥ 1026
=
207.2
= 1.082 1
E xamplE 12 Calculate the lattice constant ‘a’ of a substance having fcc lattice, molecular weight 60.2 and
density 6250 kg/m3. (N=6.02 ¥ 1026/kg mole)
Solution Given molecular weight M = 60.2, density r = 6250 kg/m3, and N = 6.02 ¥ 1026/kg mole.
For fcc lattice n = 4
1/ 3
Ê 4M ˆ
Lattice constant a = Á
Ë N r ˜¯
1/ 3
È 4 ¥ 60.2 ˘
\ a=Í 26 ˙
Î 6250 ¥ 6.02 ¥ 10 ˚
1/ 3
È 240.8 ˘
=Í 26 ˙
Î 37265 ¥ 10 ˚
= 4Å
E xamplE 13 In NaCl crystal, the spacing between the successive (1 0 0) plane is 2.82 Å. X-ray incident on
the surface of the crystal is found to give rise to first order Bragg reflection at glancing angle 8.8º. Calculate
the wavelength of X-ray.
Solution Given d = 2.82 Å, q = 8.8º and n = 1.
Formula used is 2d sin q = nl.
2 ¥ 2.82 ¥ 10–10 ¥ sin 8.8 = l
fi l = 5.64 ¥ 10-10 ¥ 0.153
= 0.863 Å
E xamplE 14 The first-order diffraction is found to occur at a glancing angle of 9º. Calculate the wavelength of
X-ray and the glancing angle for second order diffraction if the spacing between the adjacent plane is 2.51 Å.
Solution Given n =1, q = 9º and d = 2.51 Å = 2.51 ¥ 10–10 m.
Formula used, is 2d sin q = nl.
Therefore, l = 2 ¥ (2.51 ¥ 10–10) ¥ sin 9º = 0.7853 Å.
For n = 2
Ê 2l ˆ Ê lˆ
q = sin -1 Á ˜ = sin -1 Á ˜
Ë 2d ¯ Ëd¯
Ê 0.7853 ¥ 10-10 ˆ
= sin -1 Á ˜
Ë 2.51 ¥ 10-10 ¯
= 18.2
E xamplE 15 X-rays of wavelength 1.5 Å make a glancing angle 60º in the first-order when diffracted from
NaCl crystal. Find the lattice constant of NaCl.
Solution Given l = 1.5 Å, q = 60o and n = 1.
Formula used is 2d sin q = nl.
-10
Ê nl ˆ 1 ¥ 1.5 ¥ 10 1.5 ¥ 10-10 ¥ 2
\ d =Á = =
Ë 2 sin q ˜¯ 2 sin 60∞ 2¥ 3
= 0.87Å
E xamplE 16 X-ray of wavelength 1.4 Å is found to be Bragg reflected from the (1 1 1) plane of an fcc
structure. If the lattice parameter of the crystal is 5 Å. Find the angle at which the X-ray is incident on the (1
1 1) plane of the crystal.
Solution Given l = 1.4 Å, lattice parameter of fcc structure (a) = 5 Å and the plane of the fcc structure = (1 1 1).
a
Interplanar spacing d hkl =
h + k2 + l2
2
and 2d sin q = nl
a 5 ¥ 10-10
\ d111 = = = 2.887 ¥ 10-10 m
3 3
-10 ˘
Ê nl ˆ -1 È 1 ¥ (1.4 ¥ 10 )
and so, q111 = sin -1 Á = sin Í -10 ˙
Ë 2d111 ¯˜ Î 2 ¥ 2.887 ¥ 10 ˚
= 14
E xamplE 17 Calculate the glancing angle on the cube face (1 0 0) of a rock salt crystal (a = 2.184 Å)
corresponding to second order reflection of X-rays of wavelength 0.710 Å.
Solution Given d = a = 2.814 Å = 2.814 ¥ 10–10 m for cube face (1 0 0), n = 2 for second order diffraction and l = 0.710
¥ 10–10 m.
2d sin q = nl
È nl ˘
q = sin -1 Í ˙
Î 2d ˚
È 2 ¥ 0.710 ¥ 10-10 ˘
= sin -1 Í -10 ˙
Î 2 ¥ 2.814 ¥ 10 ˚
= 14.6
E xamplE 18 From the following data calculate the wavelength of neutron beam and its speed. Spacing
between successive planes is 3.84 Å, glancing angle 30° and the order of Bragg reflection = 1.
Solution Given d = 3.84 Å = 3.84 ¥ 10–10 m, q = 30° and n = 1.
h
Formula used are 2d sin q = nl and l = .
mn
Thus, 2 ¥ 3.84 ¥ 10–10 ¥ sin 30° = l
1
or l = 2 ¥ 3.84 ¥ 10-10 ¥ = 3.84Å
2
h
 l= [according to deBroglie relation]
mn
h 6.62 ¥ 10-34
\ l= =
ml 1.67 ¥ 10-27 ¥ 3.84 ¥ 10-10
l = 1.03 ¥ 103 m
E xamplE 19 Electrons accelerated from the state of rest by 120 V are reflected from an fcc crystal. The
reflection maximum is observed at 22°. Determine the lattice parameter if the Bragg reflection occurs from
the (1 1 1) plane.
Solution Given V = 120 V, q = 22° and n = 1.
Formula used are l =

h È 1 2 ˘
(2meV)1/ 2 ÍÎ 2 mv = eV ˙˚
and 2d sin q = nl
h
\ l=
2meV
6.6 ¥ 10-34
= = 1.12 ¥ 10-10 m
-31 -19
2 ¥ 9.1 ¥ 10 ¥ 1.6 ¥ 10 ¥ 120
nl 1 ¥ 1.12 ¥ 10
so, d111 = = = 1.4949Å
2 sin q 2 ¥ sin 22∞
a
and d111 =
3
\ a = d111 ¥ 3 = 3 ¥ 1.4949Å = 2.589Å
ExamplE 20 A monochromatic beam of X-rays of wavelength 1.24 Å is reflected by cubic crystal of KCl.
Determine the interplanar distances for (1 0 0), (1 1 0) and (1 1 1) planes. Given density of KCl = 1980 kg/
m3 and molecular weight M = 74.5. Avogadro’s number N = 6.023 ¥ 1026/kg mole.
Solution Given M = 74.5, r = 1980 kg/m3 and N = 6.023 ¥ 1026/kg mole.
Formulas used are
nM
a3 =
Nr
a
and d hkl = for cubic crystal.
h + k2 + l2
2
nM 4 ¥ 74.5
\ a3 = = = 24.99 ¥ 10-29
N r 6.023 ¥ 1026 ¥ 1.98 ¥ 103
a = [249.9 ¥ 10-30 ]1/ 3 = 6.30 ¥ 10-10 = 6.3Å
a a 6.3Å
d100 = = = = 6.3Å
12 + 02 + 0 1 1
a a
d110 = = = 4.38 ¥ 10-10 m = 4.38Å
2 2
1 +1 + 0 2
a a 6.3 ¥ 10-10
d111 = = = = 3.64 ¥ 10-10 m = 3.64Å
12 + 12 + 12 3 3
E xamplE 21 Determine the potential energy of K+ and Cl– ion when they are separated by a distance of
0.15 nm.
Solution Given separation distance r0 = 0.15 nm = 0.15 ¥ 10–9 m.
-e2 -e
The potential energy of the ions (V) = J= eV
4pe 0 r0 4pe 0 r0
1.6 ¥ 10-19
\ V=
4 ¥ 3.14 ¥ 8.85 ¥ 10-12 ¥ 0.15 ¥ 10-9
= 0.0959 ¥ 102 eV
= 0.096 ¥ 102 eV = – 9.6 eV
E xamplE 22 From the following data determine the cohesive energy of NaCl. The equilibrium separation
r0 = 0.32 nm, a = 1.748, n = 9, ionisation energy = 4 eV and electron affinity = –2.16 eV.
Solution Given r0 = 0.32 nm = 0.32 ¥ 10–9 m, a = 1.748, n = 9, Ionisation energy = 4 eV and electron affinity = –2.16 eV.
-a e 2 Ê 1ˆ -a e Ê 1ˆ
V (r0 ) = Á1 - ˜ Joule or V ( r0) = Á1 - ˜ electron volt
4pe r0 Ë n¯ 4pe 0 r0 Ë n¯
-1.748 ¥ (1.6 ¥ 10-19) 2 Ê 1ˆ
= Á1 - ˜
4 ¥ 3.14 ¥ 8.85 ¥ 10-12 ¥ 0.32 ¥ 10-9 Ë 9 ¯
= 0.0698 eV
E xamplE 23 Find the ratio of number of Schottky defects to the total number of cation–anion pairs for a
binary ionic crystal of NaCl of the average energy required to produce a Schottky defect be 2.02 eV at room
temperature. Given Boltzman constant k = 1.38 ¥ 10–23 J/K.
-U ˆ
Solution The number of Schottky defects is given by n = N exp ÊÁË ˜¯ where N is the total number of cation–anion
2kT
pairs.
Room temperature T = 27°C = 300 K.
n Ê -2.02 ¥ 1.6 ¥ 10-19 ˆ
Hence = exp Á ˜
N Ë 2 ¥ 1.38 ¥ 10-23 ¥ 300 ¯
= 1.12 ¥ 10-17

Q.1 Which one of the following Bravais lattices is not found in cubic crystals?
(a) simple cubic (b) body centered cubic
(c) face centered cubic (d) base centered cubic
Q.2 In hexagonal crystal, the angles between the axes are
(a) a = b = g = 90° (b) a π b π g π 90°
(c) a = b = 90°, g = 120° (d) a = g = 90° π b
Q.3 The coordination number of bcc lattice is
(a) 6 (b) 8 (c) 12 (d) 3
Q.4 The number of atoms present in the fcc lattice is

(a) 1 (b) 2 (c) 3 (d) 4
Q.5 The packing fraction of diamond crystal structure is
(a) 74% (b) 68% (c) 52% (d) 34%
Q.6 The number of ions pair per unit cell present in NaCl is
(a) 5 (b) 2 (c) 4 (d) 1
Q.7 X-rays can be deflected by
(a) magnetic fields (b) electric fields
(c) large–size obstacles (d) none of these
Q.8 The Miller indices of a set of parallel planes which make equal intercepts on the three axes are
(a) (1 2 1) (b) (1 1 1) (c) (1 0 0) (d) (1 0 1)
Q.9 In a simple cubic lattice the ratio d100 : d110 : d111 is
(a) 6 : 3: 2 (b) 3: 6 :1 (c) 6 : 3 : 2 (d) 6:3:1
Q.10 Hexagonal closed pack structure is
(a) simple cubic structure (b) body centered cubic structure
(c) face centered cubic structure (d) base centered cubic structure
Q.11 Which one of the following crystals is an example of monoclinic-
(a) NaCl (b) Mg (c) CaSO4 (d) CuSO4.5H2O
Q.12 Which one of the following crystals is the example of rhombohedral (or trigonal)-
(a) CsCl (b) As (c) Fe3C (d) K2Cr2O7
Q.13 The atomic radius for fcc lattice is
a a 3 a
(a) (b) (c) a (d)
2 2 2 4 2
Q.14 The packing fraction has maximum value for
(a) fcc structure (b) sc structure (c) hpc structure (d) bcc structure
Q.15 NaCl crystal has
(a) fcc structure (b) hpc structure (c) sc structure (d) bcc structure
Q.16 Electrons are transferred from one type of atom to the other type in
(a) covalent crystals (b) ionic crystals (c) metallic crystals (d) molecular crystals
Q.17 van der Waals bond is formed due to
(a) the sharing of valence electrons
(b) attraction between positive ion cores and electrons
(c) electrostatic attraction between oscillating or permanent dipoles
(d) electrostatic force of attraction between ions
Q.18 Which one of the following bonds has high electrical conductivity?
(a) ionic bond (b) covalent bond (c) metallic bond (d) H-bond
Q.19 The energy required to separate the atoms of a molecule is called
(a) ionisation energy (b) Coulomb energy
(c) cohesive energy (d) dissociation energy
Q.20 If n be the number of atoms in a unit cell of the cubic system, N and M be the Avogadro’s number and atomic
weight, respectively, and 3 be the density of the element, then the lattice constant a is given by
1/ 3 -1/ 3 1/ 3
È nM ˘ È nM ˘ ÈM r˘ nM
(a) Í (b) Í ˙ (c) Í (d)
˙
Î Nr ˚ Î Nr ˚ Î nN ˙˚ Nr
Q.21 Which of the following shape of unit cells are correct for most of the crystals
(a) parallelopiped (b) elliptical (c) spherical (d) none of these
Q.22 A cubic system is completely represented by
(a) a = b = c (b) a = b π c
a = b = g = 90° a = b = g = 90°
(c) a = b = c (b) none of these
a = b = g π 90°
Q.23 The Bravais lattice of CsCl structure is
(a) body centered cubic (b) face centered cubic
(c) simple cubic (d) none of these
Q.24 The arrangement of atoms in a crystal is known as
(a) crystal structure (b) lattice
(c) Bragg’s lattice (d) none of these
Q.25 Covalent bond is formed
(a) by emission of electron from the atom
(b) by transferring of electrons from one to another atom
(c) by sharing of electrons
(d) none of these
Q.26 The potential energy between atoms in equilibrium is
(a) minimum (b) maximum (c) both (a) & (b) (d) none of these
Q.27 The resultant force between the atoms in equilibrium is
(a) large (b) zero (c) attractive (d) none of these
Q.28 Inter molecular bonds are
(a) H-bonds (b) dipole bonds (c) dispersion bonds (d) all of these
Q.29 If the Miller indices of a plane is (1 0 0), then
(a) the plane is perpendicular to x-axis (b) the plane parallel to x-axis
(c) the plane is perpendicular to y-axis (d) none of these
Q.1 What is a crystal?

Q.2 Why are most solids crystalline in nature?
Q.3 Define crystal lattice, plane lattice and space lattice.
Q.4 What is crystal structure?
Q.5 Define a primitive unit cell. Can a unit cell be primitive?
Q.6 Define mathematically a lattice. What conditions an ideal lattice must follow?
Q.7 Define translation vector, unit cell, coordination number and packing efficiency.
Q.8 Write coordination number of simple cubic, body centered cubic and face centered cubic lattices.
Q.9 What is coordination number of diamond crystal? Also state the number of next nearest neighbour.
Q.10 What is coordination number of hcp structure?
Q.11 How many lattice points are there in the cubic unit cell of bcc structure.
Q.12 Explain diamond structure. Calculate its packing fraction. Give examples of any two materials having
this structure.
Q.13 Why X-rays are used for crystal structure analysis?
Q.14 Why g-rays are not used to study crystal structure?
Q.15 What is Bragg’s law?
Q.16 What is Bragg’s equation?
Q.17 How does Bragg reflection differ from ordinary reflection?
Q.18 Explain the term bonding.
Q.19 What are different kinds of bonding?
Q.20 What is ionic crystal?
Q.21 Name various point defects in solids.
P ractice P roblems
general Questions
Q.1 (a)
Distinguish between a crystal and an amorphous solid.
(b)Give three main differences between crystalline and amorphous solids.
Q.2 (a)
What is crystal structure? State the relation between crystal structure, lattice and basis.
(b)Define a primitive cell. Distinguish between a primitive unit cell and non-primitive unit cell with
the help of diagram. Can a unit cell be primitive?
Q.3 What is the concept of Miller indices? Derive the formula for the distance between two adjacent planes
of a simple cubic lattice.
Q.4 What is Bravais lattice? Explain different types of Bravais lattices in three dimensions.
Q.5 Draw the diagrams of the following structures: NaCl and CsCl. Give at least two examples of each
structure.
Q.6 Explain the crystal structure of sodium chloride (NaCl). Draw a sketch of sodium chloride lattice and write
down the coordinates of the atoms in the unit cell. What is the number of sodium ions in unit cell of NaCl?
Q.7 Explain the Crystal structure of diamond. In diamond crystal, what is the number of nearest neighbours,
the number of atoms per unit cell and packing fraction? Show that it has comparatively loose packing.
Q.8 (a) Explain the concept of Miller indices. How are they calculated? How the orientation of a plane is
specified by Miller indices? Define Miller indices of a direction. State their important features.
(b) Why the reciprocals of intercepts of the plane are taken to find Miller indices?
– –
Q.9 Draw the planes (1 0 0), (0 1 0), (0 0 1), (1 1 0), (1 0 1), (0 1 1), (2 0 0), (2 0 0), (1 0 0), (2 0 1), (1 1 1),
and (1 1 2) in a simple cubic unit cell.
Q.10 Derive the expression for the interplanar spacing between two parallel planes with Miller indices
(h k l) and show that for a simple cubic lattice of lattice constant a
a
d hkl =
h2 + k 2 + l 2
Q.11 Derive Bragg’s law of crystal diffraction 2d sin q = nl and give its significance. Discuss briefly the
method of crystal structure determination.
Q.12 Is there any interdependence of coordination number and packing efficiency. Illustrate by giving
examples.
Q.13 Discuss briefly the experimental method for crystal structure determination by X-ray diffraction.
Q.14 Name the standard experimental methods of X-ray diffraction.
Q.15 Describe in detail Laue method and also describe the usefulness of this method.
Q.16 Explain with necessary theory the powder method for X-ray analysis.
Q.17 Describe in detail the powder method and its usefulness.
Q.18 What are point defects in solids?
Q.19 What are different types of point defects? Explain.
Q.20 What are Schottky and Frenkel defects. Derive the necessary relation to show that Schottky defects in
ionic crystal depend on the temperature.
Q.21 Show that the number of Frenkel defects in equilibrium at a given temperature is proportional to
(NNi)1/2, where N and Ni are number of atoms and interstitial atoms respectively.
Q.22 Name various types of bonds in solids and given one example of each.
Q.23 Explain any four types of bondings in solids.
Q.24 Write short note on bonding in solids.
  
Q.1 Find the Miller indices for planes in each of the following sets which intercept a , b , and c , at (i) 3a,
3b, 2c and (ii) a, 2b, •. [Ans: (2,2,3) and (2,1,0)]
Q.2 Lattice constant of a cubic lattice is a. Calculate the spacing between (i) (011), (ii) (101) and, (iii) (110)
planes. [Ans: (i) a / 2 (ii) a / 2 (iii) a / 2 ]
Q.3 For a cubic lattice calculate the distance of (123) and (234) planes from a plane passing through the origin.
[Ans: a / 14 and a / 29 ]
Q.4 Calculate the glancing angle at which X-rays with l =1.549 Å will be reflected in first and second
orders from a crystal with interplanar distance 4.225 Å. [Ans: q1 = 10°31¢ and q2 = 21°21¢]
Q.5 Using 2.02 as the value of lattice constant, calculate the wavelength of X-ray in second order, if angle
of diffraction q = 26°. [Ans: 1.24 Å]
Q.6 A crystal is mounted on an X-ray spectrometer. X-rays are incident at the glancing angle for three
reflections are 5°28¢, 12°1¢ and 18°12¢. Show that these are successive orders of reflections from the
same crystal plane. Also find the spacing. [Given l for X-rays used as 0.586 Å].
[Ans: 2.817 Å, 2.817 Å and 2.817 Å]
Q.7 A certain crystal reflects monochromatic X-rays strongly when Bragg glancing angle (first order) is
15°. What are the glancing angles for second and third order spectrum. [Ans: 31.17°, 50.93°]
Development of Quantum
Mechanics
15
Learning Objectives
LO 1 Learn about blackbody radiation and LO 4 Discuss the de Brogile waves and its
Planck’s quantum hypothesis demonstration by Davisson-Germer
LO 2 Understand the concept of quantum experiment
theory LO 5 Explain Compton effect and its
LO 3 Know about wave particle duality and verification
photoelectric effect and its theoretical LO 6 Evaluate phase and group velocities
applications and their interrelationship
Introduction
Newton’s laws describe the motion of particles in classical mechanics and Maxwell’s equations describe
the electromagnetic fields in classical electromagnetism. The classical mechanics correctly explains the
motion of celestial bodies like planets, stars, macroscopic and microscopic terrestrial bodies moving
with non-relativistic speeds. However, classical theory does not hold in the region of atomic dimensions,
i.e., it cannot explain the non-relativistic motion of electrons, protons etc. Classical theory could not
explain the stability of atoms, spectral distribution of blackbody radiation, the origin of discrete spectra
of atoms, etc. Also, classical mechanics could not explain a large number of observed phenomena like
photoelectric effect, Compton effect, Raman effect, etc. So, the insufficiency of classical mechanics led to
the development of quantum mechanics. Quantum mechanics is the description of motion and interaction
of particles at the small scales where the discrete nature of the physical world becomes important. The
quantum mechanics for the atomic system led to the explanation of discrete energy levels as well as
the postulation of different quantum numbers. Niels Bohr had a large influence on the development of
quantum mechanics through his so-called Copenhagen Interpretation, a philosophical construct that was
formulated to provide a fundamental framework for understanding the implicit assumptions, limitations,
and applicability of the theory of quantum mechanics.
The development of quantum mechanics took place in two stages. The first stage began with Max Planck’s
hypothesis according to which the radiation is emitted or absorbed by matter in discrete packets or quanta
of energy. This energy is equal to hn, where h is Planck’s constant and n is the frequency of radiation.
This hypothesis led to a theory which was not completely satisfactory being a mixture of classical and
non-classical concepts. The second stage of quantum mechanics began in 1925 along with two points of
views. For example, matrix mechanics was introduced by Heisenberg, in which only observed quantities
like frequencies and intensities of spectral lines are taken into account and unobserved quantities like
positions, velocities, etc. in electronic orbits are omitted. Another form of quantum mechanics is called
wave mechanics, whose theory was developed by Schroedinger in 1926. In this mechanics, concepts of
classical wave theory and deBroglie’s wave particle relationship are combined with each other. With the
application of quantum mechanics, several problems of atomic physics have been solved. However, this
mechanics also has certain limitations. Therefore, a more complete theory of particles called quantum
field theory has been accepted since 1947. In order to understand the development of wave mechanics,
we begin with the blackbody radiation.
15.1 BLackBOdy RadiatiOn: SpectRaL diStRiButiOn LO1

A body that completely absorbs radiations of all wavelengths incident on it is referred to as a blackbody.
When such a body is heated, it emits radiations which we call as blackbody radiations. A cavity made out of
a hollow container of any material (iron or copper) with
a narrow opening and painted with lampblack in the
inside portion gives a close approximation to a perfectly 2000 K
blackbody. When any radiation falls on this hole, it
1750 K
enters the cavity, gets reflected by the wall of the cavity 1000 K
Intensity I(ν )
and eventually gets absorbed. Now if we heat the 1250 K

container at various temperatures, it will emit radiations
of all the frequencies (or wavelengths). So, the emitted
radiation from a blackbody is a continuous spectrum,
1500 K
i.e., it contains radiation of all the frequencies.
Let the intensity of emitted radiation be I(n) dn
between the frequencies n and n+dn. The experimental 1 2 3 4 5 6 7 8
measurements of intensity I(n) with different n is shown ν ( ×1014)
in Fig. 15.1 for different values of temperature T. These
Figure 15.1
plots show that
(1) The distribution of frequencies is a function of temperature of the blackbody.
(2) With the increase in temperature, the total amount of emitted radiation I(n) dn increases.
(3) The position of the maximum peak shifts toward higher frequencies with increasing equilibrium
temperature.
The classical electromagnetic theory or wave theory together with classical thermodynamics does not explain
the characteristics of the blackbody spectrum. However, Planck’s hypothesis can explain these characteristics
together with the use of classical thermodynamics.
Development of Quantum Mechanics 553
Classical wave theory says that the electromagnetic radiation inside the cavity of the blackbody at an
equilibrium temperature T forms the standing waves and the number of standing waves (possible modes) that
can fit in the cavity depends on the wavelength. The number of possible modes in the cavity is large if the
wavelength is small. However, for large wavelengths the number of possible modes is small. According to
Rayleigh and Jeans, this increase in the number of modes is proportional to 1/ l2 or n2 and also each of the
standing waves must be assigned an average kinetic energy kT, where k is the Boltzmann constant. This leads
to the following Rayleigh-Jeans law (details discussed later)
Rayleigh-Jeans Law
8pn 2
I (n )dn = kTdn
c3
This relation shows that I(n) is proportional to the square of n. Experimental
I(ν )
Result
The corresponding plot is shown in Fig. 15.2. It is clear from the
figure that the experimental data does not agree with the theory; the
agreement is good only for smaller values of n. The disagreement
at high frequencies, i.e., in the UV region, is called ultraviolet
catastrophe. Thus, the spectral distribution of a blackbody could ν
not be explained on the basis of classical theory. This difficulty
Figure 15.2
was resolved by Planck in 1900, when he stated that by assuming
electromagnetic radiation to be emitted or absorbed in bundles of size hn, one could correctly predict the
spectrum of blackbody radiation. As mentioned earlier, this bundle of energy is called a quantum. The quanta
of high frequencies have high energies and those of low frequencies have low energies. Thus, the atoms and
molecules in the cavity will emit radiation only if they have energy in the excess of hn. For low frequencies
n, there will be a large number of atoms and molecules that might have this excess energy. Since the bundles
become quite bigger for higher frequencies n, the number of atoms or molecules having energies in the excess
of hn decreases. It means for large n, the intensity I(n) does not increase rather decreases.
For the explanation of blackbody radiation, Planck made a use of the Maxwell-Boltzmann distribution.
According to this distribution, the number of molecules Nn with energy E is given by
Nn = N0e–E/kT
In the above expression, N0 refers to the number of molecules

in a system in equilibrium at temperature T. Planck combined
the expression of Nn with his quantum hypothesis E = nhn and
calculated the mean energy. Finally, he arrived at the following
expression for the distribution of the maximum intensity of
radiation in the spectrum of blackbody.
I(ν )
8p hn 3 1
I (n ) = 3 ( hn / kT )
c e -1
This expression is referred to as Planck’s radiation law. This
theoretical formula fits very well with the experimental data for 1 2 3 4 5 6 7 8
the entire wavelength, as shown in Fig. 15.3. Thus, Planck’s ν ( ×1014)
quantum theory was able to interpret fully different characteristics Figure 15.3
of blackbody radiation which classical theory could not.
15.2 pLanck’S Quantum HypOtHeSiS LO1

Max Planck in 1900 introduced the quantum theory of radiation to explain the distribution of energy in the
spectrum of blackbody radiation. He assumed that the atoms of the walls of a blackbody behave as oscillators
and each has a characteristic frequency of oscillation. He made the following two revolutionary assumptions
about the atomic oscillator.
(1) An oscillator cannot have any arbitrary value of energy but can have only discrete energies as per
the following relation
E = nhn
where n = 0,1,2,…; n and h are known as frequency of oscillation and Planck’s constant
(= 6.62 ¥ 10–34 J sec), respectively. This relation shows that the energy of the oscillation is quantised.
(2) The oscillator can emit or absorb energy only in the form of packets of energy (hn) but not
continuously. In other words, we can say that the emission or absorption of energy occurs only when
the oscillator jumps from one energy state to another along with the energy difference given by
DE = Dnhn
or E2 – E1 = (n2 – n1)hn
15.2.1 Average Energy of Planck’s Oscillators

If N be the total number of oscillators and E as the total energy of these oscillators, then the average energy
will be given by the relation
E (i)
E=
N
Now consider N0, N1, N2,…Nn, as the number of oscillators having energy values 0, hn, 2hn, …, nhn,
respectively. Then by Maxwell’s distribution formula Nn = N0e–nhv/kT, we have
N = N 0 + N1 + N 2 + ....
= N 0 (1 + e - hn / kT + e -2 hn / kT + ....)
N0
= (ii)
(1 - e - hn / kT )
and the total energy is

E = ( N 0 ¥ 0) + ( N1 ¥ hn ) + ( N 2 ¥ 2hn ) + ....
= ( N 0 ¥ 0) + ( N 0 e - hn / kT ¥ hn ) + ( N 0 e -2 hn / kT ¥ 2hn ) + ....
= N 0 e - hn / kT hn[1 + 2e - hn / kT + 3e -2 hn / kT + ...]
hn
= N 0 e - hn / kT - hn / kT 2 (iii)
(1 - e )
Putting the values of N and E from above relation in Eq. (i), we get
E hn e - hn / kT hn
E= = - hn / kT
= hn / kT
N 1- e e -1
hn
E= (iv)
e hn / kT - 1
This is the expression for average energy of a Planck’s oscillator.
15.2.2 Planck’s Radiation Formula

The energy density of radiation (un) in the frequency range n and n + dn depending upon the average energy
of an oscillator is given by
8pn 2
un dn = dn ¥ E (v)
c3
Putting the value of E from Eq. (iv) in Eq. (v) gives
8pn 2 hn
un dn = dn
c3 e hn / kT - 1
8p h n 3 dn
or un dn = (vi)
c 3 e hn / kT - 1
The above relation is known as Planck’s radiation formula in terms of frequency. This law can also be ex-
c c
pressed in terms of wavelength l of the radiation. Since n = for electromagnetic radiation, dn = - 2 d l.
l l
Further, we know that the frequency is reciprocal of wavelength or in other words an increase in frequency
corresponds to a decrease in wavelength. Therefore
uldl = –undn
3
Ê cˆ Ê c ˆ
Á ˜ - d l˜
8p h Ë l ¯ ÁË l 2 ¯
or ul d l = - 3 hc / l kT
c e -1
8p hc 1
ul d l = 5 hc / l kT
dl (vii)
l e -1
The above relation is known as Planck’s formula in terms of wavelength.
15.2.3 Wien’s Law and Rayleigh-Jeans Law
With the help of Planck’s radiation formula Wien’s law and Rayleigh-Jeans law can be derived. When the
wavelength l and temperature T are very small, then ehc/lkT >> 1. Therefore, 1 can be neglected in the
denominator of Eq. (vii).
Thus
8p hc
ul d l = e - hc / l kT d l
l5
hc
By substituting 8phc = A and = B, we get
k
A
ul d l = e - B / lT d l (viii)
l5
This is known as Wien’s law, which is valid at low temperature T and small wavelength l.
hc
For high temperature T and large wavelength l, ehc/lkT can be approximated to 1 + . Then we have from
Eq. (vii) l kT
8p hc
ul d l = dl
5Ê hc ˆ
l Á1 + - 1˜
Ë l kT ¯
8p kT (ix)
ul d l = dl
l4
This is known as Rayleigh-Jeans Law.
15.3 SimPLE COnCEPt OF QuAntum thEORy LO2

As discussed earlier, Planck’s hypothesis says that the radiation does not emit in continuous fashion rather
it gets emitted in discrete packets of energy equal to hn. These packets are referred to as quanta or photons.
Therefore, it can be said that the exchange of energy between the radiation and the matter takes place in
discrete set of values. In view of the application of quantum theory, it is necessary to be aware of the photon.
15.3.1 Photon: mass, Energy and momentum

Photon is an elementary particle that is massless and has no charge. It is a bundle of energy or packet of
energy emitted by a source of radiation. It moves with velocity of light. It can carry energy and momentum.
We know that the mass m of the particle moving with v, comparable with the velocity of light c, is given by
as per the special theory of relativity.
m0
m= (i)
1 - v2 / c2
where m is the relativistic mass of the particle and m0 is its rest mass. Since the photon is moving with the
velocity of the light, we substitute v = c in Eq. (i). With this the moving mass m of the photon becomes
m = •, which is not possible. So, if the photon moves with the velocity of light then the zero in the numerator
balances the zero in the denominator, i.e., m = 0/0; this is an indeterminate quantity. It means if we take rest
mass of the photon to be zero, this value should not particularly disturb us due to the fact that the photons are
never at rest and always keep moving with the velocity of light.
The energy of a photon is given below as
E = hn (ii)
If m is the moving mass of the photon, then according to special theory of relativity, the following relation
gives the energy
E = mc2 (iii)
so mc2 = hn
hn (iv)
or m=
c2
Now the energy relation
E2 = p2c2 + m02 c4 (v)
Since m0 = 0, E = pc and the momentum of the photon is given by
E mc 2
p= = = mc (vi)
c c
E hn (vii)
or p= =
c c
Thus, if a photon of frequency n is to be treated as a particle, then the characteristics of the photon are given as
m0 = 0, E = hn, m = hn/c2 and p = hn/c (viii)
These characteristics of the photons are useful in the discussion of Compton effect, which establishes the
photon hypothesis.
15.4 Wave paRticLe duaLity LO3

The phenomena of interference, diffraction and polarisation can be explained on the basis of wave theory
of light. However, the wave nature of light fails to explain the phenomena of Compton effect, photoelectric
effect, the continuous X-ray spectrum and the blackbody radiation. In the light of these facts, physicists
assumed the particle nature of electromagnetic radiation (light). These entire phenomena can be explained on
the basis of quantum hypothesis, according to which electromagnetic radiation is propagated in small packets
or bundles. These packets are called photons. It means that light or electromagnetic radiation exhibits wave
and particle properties both. Hence, light or electromagnetic radiation has dual nature, i.e., it behaves like a
particle as well as a wave. This dual characteristic property of radiation is called dual nature of light or wave
particle duality.
15.5 PhOtOELECtRiC EFFECt LO3

The photoelectric effect refers to the emission or ejection of electrons from the surface of a metal (generally)
in response to incident light. Energy contained within the incident light is absorbed by the electrons within the
metal, gaining sufficient energy to be ‘knocked’ out of, i.e., emitted from the surface of the metal. According
to the classical Maxwell wave theory of light, the more intense incident light should eject the electrons from
the metal with their greater energy. It means the average energy carried by an ejected (photoelectric) electron
should increase with the intensity of the incident light. In fact, Lenard found that this was not so. Rather, he
observed the energies of the emitted electrons to be independent of the intensity of the incident radiation. In
1905, Einstein resolved this paradox successfully by proposing that the incident light consists of individual
quanta, called photons, that interact with the electrons in the metal like discrete particles, rather than as
continuous waves. He adopted the Planck’s quantum hypothesis and applied it to the electromagnetic radiation.
For a given frequency (n) or colour (l) of the incident radiation, each photon carries the energy E = hn.
According to Einstein’s model, increase in the intensity of the light corresponds to the enhancement in the
number of incident photons per unit time (flux), while the energy of each photon remains the same as long
as the frequency of the radiation is kept constant. It means, increasing the intensity of the incident radiation
would cause greater numbers of electrons to be ejected and each electron would carry the same average
energy because each incident photon carries the same energy. Likewise, in Einstein’s model, increasing the
frequency n rather than the intensity of the incident radiation would increase the average energy of the
emitted electrons. Both of these predictions were confirmed experimentally. It is interesting to note that
the rate of increase of the energy of the ejected electrons with increasing frequency enables us to determine the
value of Planck’s constant h, as the frequency can be measured.
15.5.1 theoretical Explanation
In photoemission, one quantum is absorbed by one electron. If the electron is some distance into material of
the cathode, some energy will be lost as it moves towards the surface. There will always be some electrostatic
cost as the electron leaves the surface. This is known as the work function f0. The electrons those are very
close to the surface will be the most energetic, and they will leave the cathode with kinetic energy given by
EK = hn – f0
or EK = hn – hn0
where hn0 = f0
Therefore, it is clear that there is a minimum light frequency called threshold frequency n0 for a given metal
for which the quantum of energy is equal to the work function. Light below that frequency, no matter how
bright, will not cause photoemission.
15.5.2 Experiment
Light
An experimental arrangement to the photoelectric effect is
Quartz Window
shown in Fig. 15.4. It consists of a vacuum tube A, which A
contains a metallic plate B and a charge collecting plate C. ie
When light is incident on the plate B through the quartz B C
window, electrons are ejected from the metallic surface. G
The collector is kept at positive potential V with respect V
to the metallic plate, which is at zero potential. So, due to
this positive potential the collector C collects these ejected
electrons. Therefore, a current ie is produced, which can
be measured by the galvanometer G. We can increase the
current ie by increasing the potential V until ie reaches a Battery (Polarity Reversing)
constant value, i.e., it approaches a saturation.
Figure 15.4
By using the reversing switch, we apply the negative potential
to the collecting plate C. Under this situation, the electrons are repelled by C and only those electrons whose
energy is greater than the potential energy eV will be able to reach the collector C. So we get some current in the
galvanometer G. The applied potential for which the current ie becomes zero, i.e., ie= 0, is called stopping potential
V0. The relation between the maximum kinetic energy of the electrons EK and stopping potential V0 is given as
below.
1 2
Ek = mn max = e |V0|
2
We can obtain the following results by performing detailed experiment under various conditions.
(1) The photoelectric current ie increases with the increasing intensity I of the incident radiation, if the
frequency is kept constant.
(2) There is no time lag between illumination of the metal surface and the emission of electrons.
(3) If the frequency of the incident radiation is greater than the threshold frequency n0 (certain minimum
frequency), only then the emission of electrons takes place.
(4) The maximum kinetic energy EK of the photoelectrons is independent of the intensity I of the
incident light. This is shown in Fig. 15.5 in which we observe that the stopping potential is same for
the light of three different intensities having same frequency.
(5) The maximum kinetic energy of the photoelectrons depends on the frequency of the incident
radiation. From Fig. 15.6, we observe that at different frequencies, stopping potential is also different
but the saturation current remains the same.
Ie
Ie
I1
I2 n1> n2> n3
n1
I3
n2
n3
I1 > I2 > I3
– V0 O V + – V01 V02 V03 V +
(6) There is a linear relation between maximum kinetic

energy and frequency. We show this in Fig. 15.7
for three different metals cesium, potassium and
Cs K
tungsten, which satisfy the following relation W
Ek
EK = a1n + a2 q
O
Here a1 is the slope of the straight line and a2 is the n
intercept. From the figure it is clear that though a1 remains A
the same for all surfaces, a2 is different for different metals. B
2 OA = a
The photoelectric effect is perhaps the most direct and
tan q = a1
convincing evidence of the existence of photons and the C
‘corpuscular’ nature of light and electromagnetic radiation. Figure 15.7
That is, it provides undeniable evidence of the quantisation
of the electromagnetic field and the limitations of the classical field equations of Maxwell. Albert
Einstein received the Nobel Prize in Physics in 1921 for explaining the photoelectric effect and for his
contributions to the theoretical physics.
15.6 de BROgLie WaveS: matteR WaveS

In 1924, Louis deBroglie proposed in his doctoral dissertation that there was a fundamental relation between
waves and particles. Therefore, the energy of the photon according to special theory of relativity is given by
E = hn (i)
and momentum p is
hn h (ii)
p= =
c l
Here, it can be noted that E and p are the characteristics of the partricles, and n and l are the characteristics
of the waves. From above relations, we see that these sets of quantities are related to each other by the
Planck’s constant h. deBroglie also suggested that the dual nature of electromagnetic radiation may be
extended to material particles such as electrons, protons, neutrons etc. It means that a moving particle,
whatever its nature be, has wave properties associated with it. The waves associated with these particles
are known as matter waves or deBroglie waves. The difference between the electromagnetic radiation and
elementary particles is that in the case of photons, m0 = 0 and v = c but in the case of material particles m0 π 0
and v < c. deBroglie gave the following hypothesis which is applicable to all matters, radiation and particles.
(1) If there is a particle of momentum p, its motion is associated with a wave of wavelength
h
l= (iii)
p
(2) If there is a wave of wavelength l, the square of the amplitude of the wave at any point in space is
proportional to the probability of observing, at that point in space, a particle of momentum
h
p= (iv)
l
The dual nature of matter can be proved if we could show that a beam of particles also exhibits the
phenomenon of diffraction pattern just like the electromagnetic waves show the phenomena of diffraction
and interference.
15.6.1 Demonstration of matter Waves: Davisson-Germer

Experiment on Electron Diffraction
The Davisson-Germer experiment was con-
ducted in 1927 which confirmed the deBro- Electron Gun
glie hypothesis, according to which parti-
cles of matter (such as electrons) have wave
properties. This demonstration of the wave
particle duality was important historically Detector Electron
in the establishment of quantum mechan- Electron Beam
Beam
ics and of the Schröedinger equation. The (54 Volt)
experimental setup is shown in Fig. 15.8(a).
Here the electrons from a heated filament or 50°
electron gun were accelerated by a voltage f
V and allowed to fall on the surface of nickel
target. Davisson and Germer measured the
Nickel
intensity of the scattered electrons as a func- Nickel Target Target
tion of the angle f and plotted it in the form
Figure 15.8(a) Figure 15.8(b)
of polar diagram. Fig. 15.8(b) shows the
results from the accelerating voltage of 54 V. For this case, there is an intense scattering or a pronounced
peak at an angle of q = 50º. Such deflection can be ex- Incident Diffracted

plained by assuming that the electron beam has a wave Beam Beam
a
associated with it. This situation is similar to the Bragg’s
diffraction. So the waves associated with the electron
beam were satisfying Bragg’s law, which caused a dif-
fraction peak. In order to prove this, consider Fig. 15.9
d
that shows atomic planes in Ni crystal. Here q = 50º, f
q
f = (180 – 50)/2 = 65º, d = 0.91 Å. Hence, for n = 1,
q
Bragg’s law nl = 2d sin q gives a
l = 2d sin q = 2 ¥ 0.91 ¥ sin 65º = 1.65 Å (i) d

f
Since Bragg’s law basically talks about the diffraction of 2
X rays, this experiment enables us to treat the electrons as
waves and the wavelength associated with the electrons Figure 15.9
should be 1.65 Å, if they are scattered at q = 50°.
Now, we apply deBroglie’s hypothesis. Since the electron of mass m gains the velocity v when it gets
accelerated through a potential difference of V, we write the following relation for the energy for the
nonrelativistic motion of the electron
1 2
mv = eV
2
So the deBroglie wavelength associated with the electron is given by
h h h
l= = =
p mv 2meV
150
or l= Å
V
Therefore, deBroglie wavelength associated with the electron that is accelerated by 54 V is given as
150 150
l= Å= Å = 1.67 Å (ii)
V 54
A comparison of Eq. (i) with Eq. (ii) shows that the value of the wavelength l is the same in both the cases.
It means there is a wave called deBroglie wave associated with the electrons. Therefore, this confirms the
deBroglie hypothesis.
15.7 COmPtOn EFFECt: COmPtOn SCAttERinG LO5

As per classical electromagnetic theory, when an electromagnetic radiation (frequency n) is incident on free
charges (say, electrons), the free charges absorb this radiation and start oscillating at frequency n. Then these
oscillating charges radiate electromagnetic waves of the same frequency n. This type of scattering where the
change in frequency or wavelength does not take place is called coherent scattering. This coherent scattering
has been observed with the radiation in visible range and also at longer wavelengths. However, this predica-
tion of classical theory fails in the case of scattering of radiation of very short wavelengths like X-rays. Here
the scattered X-rays are found to consist of two frequencies: n and n¢. The wavelength l corresponding to the
frequency n is called unmodified wavelength or unmodified radiation, whereas the wavelength l¢ correspond-
ing to the frequency n¢ is called modified wavelength or modified radiation. This type of scattering is known
as incoherent scattering.
The Compton effect or Compton scattering is related to the scattering of X-rays (electromagnetic waves of
very short wavelength) by free electrons. A.H. Compton found that when X-rays are scattered by a solid
material (say carbon in which the loosely bound electrons are assumed to be almost free) the scattered X-ray
radiations carry the longer wavelength. This h/��sin �
phenomenon of increase in the wavelength (or
decrease in frequency) of X-ray radiations by ton
ho �
scattering is called the Compton effect. This ed
P h/�
er � =
effect was explained by using the quantum att ��)
, p
Incident Photon Sc (h
� E = h/��cos �
theory of radiations. On the basis of this
theory, these radiations are made up of photons h� Target R � m 0v
ec cos �
Electron m v �1 – v2/c 2
of energy hn. These photons in the incident o
at Rest Ek = iled p�= 0
X-rays collide with the free electrons of the mc Ele �1 – v2/c 2

2
– m ctro
target (Fig. 15.10). If the collision is elastic, 0c 2
n
then the energy or wavelength of the scattered
photons remains the same. If the collision is
inelastic, then the incident photon transfers m0v
sin �
some of the energy to the electron. Thus, as a �1 – v2/c 2
result, the energy of the scattered X-ray photon Figure 15.10
decreases (or wavelength increases).
Let us assume,
l = wavelength of incident X-rays
l¢ = wavelength of scattered X-rays
Dl = l¢ – l = Compton shift
Energy of incident X-ray photon = hn = hc/l
Momentum of incident X-ray photon = h/l
Energy of the scattered X-ray photon hn¢ = hc/l¢
Momentum of the scattered X-ray photon = h/l¢
È m0 ˘
Kinetic energy of the recoiled electron = (m - m0) c 2 = Í - m0 ˙ c 2
2 2
ÎÍ 1 - v /c ˚˙
m0 v
Momentum of the recoiled electron = mv =
1 - v 2 /c 2
m = moving mass or relativistic mass of the electron and m0 is the rest mass.
According to the law of conservation of energy
Energy of the incident photon = Energy of the scattered photon + Energy of the recoiled electron.
E = E¢ + EK
hc hc
= + (m - m0 )c 2
l l¢
h h m0 c
or - = - m0 c
l l¢ 1 - v 2 /c 2
h h m0 c
or - + m0 c = (i)
l l¢ 1 - v 2 /c 2
According to the law of conservation of momentum,
h m0v
sin q = sin f (ii)
l¢ 1 - v 2 /c 2
h h m0v
and = cos q + cos f
l l¢ 1 - v 2 /c 2
h h m0v
or - cos q = cos f (iii)
l l¢ 1 - v 2 /c 2
Squaring and adding Eqs. (ii) and (iii), we get
h2 h2 2h 2 m02 v 2
+ - cos q =
l ¢2 l2 ll ¢ 1 - v 2 /c 2
h2 h2 2h 2 m2 v 2 c 2
or + - cos q = 20 2 (iv)
l ¢2 l2 ll ¢ c -v
On squaring Eq. (i), we get
h2 h2 2h 2 2hm0 c 2hm0 c m02 c 2 m02 c 4
+ + m02 c 2 - - + = =
l2 l ¢2 ll ¢ l¢ l 1 - v 2 /c 2 c 2 - v 2
h2 h2 2h 2 Ê1 1ˆ m2 c 4
or + - + 2hm0 c Á - ˜ = 2 0 2 - m02 c 2
l2 l ¢2 ll ¢ Ë l l¢¯ c - v
2 4 2 4 2 2 2
h2 h2 2h 2 Ê l ¢ - l ˆ m0 c - m0 c + m0 v c
or + - + 2hm0 c Á ˜ =
l2 l ¢2 ll ¢ Ë ll ¢ ¯ c2 - v2
2 2 2
h2 h2 2h 2 Ê l ¢ - l ˆ m0 v c (v)
+ - + 2hm0 c Á ˜ = 2
l2 l ¢2 ll ¢ Ë ll ¢ ¯ c - v2
On comparing Eqs. (iv) and (v), we get
h2 h2 2h 2 Ê l¢ - l ˆ h
2
h2 2h 2
+ - + 2hm0 c Á ˜ = 2 + 2 - cos q
l2 l ¢2 ll ¢ Ë ll ¢ ¯ l l¢ ll ¢
2
Ê l ¢ - l ˆ 2h
or 2hm0 c Á = (1 - cos q )
Ë ll ¢ ˜¯ ll ¢
h
or (l ¢ - l ) = Dl = (1 - cos q )
m0 c
h
or Dl = (1 - cos q ) (vi)
m0 c
where, h = Planck’s constant. m0 = rest mass of electron, c = velocity of light and q = angle of scattering of the
photon. This is just to emphasise that the RHS contains the angle of scattering of photon not of the electron.
15.7.1 Verification of Compton Effect

A beam of monochromatic X-ray of known wavelength is made incident on a graphite target. The intensity
distribution with wavelength of monochromatic X-ray scattered at different angles is measured by Bragg’s
X-ray spectrometer. The intensity distribution with wavelength for different angles is shown in Fig. 15.11.
It may be noted that diffraction patterns have two diffraction peaks-one corresponding to modified radiation
and the other corresponding to unmodified radiation. The difference between two peaks on the wavelength
axis provides the Compton shift. It can be concluded from the diffraction patterns that greater scattering angle
yields greater Compton shift. For example, Dl = l¢–l = h/m0c (1 – cos q) at q = 90° gives Dl = 0.024 Å.
D
l
q = 0° q = 45°
Intensity
Intensity
Wavelength l Wavelength l
Dl q = 90°
l q = 135°
∇
Intensity
Intensity
Wavelength l Wavelength l
Figure 15.11
Hence, Compton effect is experimentally verified.
15.7.2 Why Compton Effect is not observable with Visible Light?

As mentioned, the Compton effect is observed significantly with the X-rays which are very short wavelength
radiations. This can be confirmed, if we use the visible light (l = 4000 Å – 7000 Å) in place of X-rays and
calculate the Compton shift. For this we use
h
Dl = (1 - cos q ) = 0.0242(1 - cos q )Å
m0 c
For maximum shift, q = 180º. Hence
(Dl)max = 0.0484 Å
The percentage Compton shift for l = 4000 Å
( Dl ) max
= ¥ 100  0.001%
l
Similarly, the percentage Compton shift for larger wavelength of visible light, i.e., for l = 7000 Å, would be
0.0007%. So, you can see that the Compton shift for the case of visible light is not significant. For this reason,
the X-rays are appropriate for realising the Compton effect or Compton scattering.
15.8 pHaSe and gROup veLOcitieS: de BROgLie WaveS LO6
Phase Velocity
Waves have already been discussed in Chapter 1. However, here we will discuss phase and group velocities in
the context of deBroglie waves. We can write the deBroglie wave travelling along the +x direction as
y = a sin (wt – kx) (i)
where a is the amplitude, w (=2pn) is the angular frequency and k (=2p/l) is the propagation constant of the
wave. By the definition, the ratio of angular frequency w to the propagation constant k is the phase (or wave)
velocity. If we represent the phase velocity by u, then
w
u=
k
(wt – kx) is called the phase of the wave motion. It means the particle of the constant phase travels such that
(wt – kx) = constant.
d
or (w t - kx) = 0
dt
dx
w-k =0
dt
dx w
or =u= (ii)
dt k
dx
where u = is the phase (or wave) velocity. Thus the wave velocity is the velocity of planes of constant
dt
phase which advances through the medium. We can write the phase velocity u = nl and for an electromagnetic
wave E = hn, or n = E/h
h h
According to deBroglie l = =
p mn
E h mc 2 c 2
u = nl = ¥ = =
h mn mn n
c2
u= (iii)
n
Since c >> v, Eq. (iii) implies that the phase velocity of deBroglie wave associated with the particle moving
with velocity v is greater than c, the velocity of light.
Group Velocity
As we have seen, the phase velocity of a wave associated with a particle comes out to be greater than the
velocity of light. This difficulty can be overcome by assuming that each moving particle is associated with a
group of waves or a wave packet rather than a single

wave. In this context, deBroglie waves are represented
by a wave packet and hence we have group velocity
associated with them. In order to realize the concept
of group velocity, we consider the combination of two
waves, resultant of which is shown in Fig. 15.12. The
Figure 15.12
two waves are represented by the following relations
y1 = a sin (w1t – k1x) (i)
and y2 = a sin (w2t – k2x) (ii)
Their superposition gives
y =y1 + y2 = a [sin (w1t – k1x) + sin (w2t – k2x)]
È (w + w 2 )t (k1 + k2 ) x ˘ È (w1 - w 2 )t (k1 - k2 ) x ˘

or y = 2a sin Í 1 - ˙˚ cos ÍÎ - ˙˚
Î 2 2 2 2
È (w - w 2 )t (k1 - k2 ) x ˘ (iii)
\ y = 2a cos Í 1 - ˙˚ sin (w t - kx)
Î 2 2
w1 + w 2 k +k
where w= ,k = 1 2
2 2
Eq. (iii) can be re-written as
È ( D w )t ( D k ) x ˘
y = 2a cos Í - sin (w t - kx ) (iv)
Î 2 2 ˙˚
where Dw = w1 – w2 and Dk = k1 – k2.
The resultant wave Eq. (iv) has two parts.
(i) A wave of frequency w, propagation constant k and the velocity u, given by
w 2pn
u= = = nl
k 2p /l
which is the phase velocity or wave velocity.
Dw
(ii) Another wave of frequency Dw/2, propagation constant Dk/2 and the velocity G = . This
Dk
velocity is the velocity of envelope of the group of waves, i.e., it is the velocity of the wave packet
(shown by dotted lines) and is known as group velocity.
For the waves having small difference in their frequencies and wave numbers, we can write
Dw ∂w ∂(2pn ) ∂n ∂n
G= = = = = - l2
Dk ∂k ∂(2p /l ) ∂(1/l ) ∂l
l 2 ∂w
G=-
2p ∂l
This is the expression for the group velocity.
15.8.1 Relation between Group Velocity and Phase Velocity

If u be the phase (wave) velocity, then the group velocity can be written as
dw d È w˘
G= = (uk ) Íu = ˙
dk dk Î k˚
du
or G=u+k
dk
2p 2p k l
But k= fi dk = - 2
d l and =-
l l dk dl
Therefore, the group velocity is given by
Ê l ˆ
G = u + Á-
Ë d l ˜¯
du
du
or G=u-l
dl
This relation shows that the group velocity G is less than the phase velocity u in a dispersive medium where
u is a function of k or l. However, in a non-dispersive medium, the velocity u is independent of k, i.e., the
du
wave of all wavelength travel with the same speed, i.e., = 0. Then G = u. This is true for electromagnetic
dl
waves in vacuum and the elastic waves in homogenous medium.
15.8.2 Relation between Group Velocity and Particle Velocity

Consider a material particle of rest mass m0. Let its mass be m when it moves with velocity v. Then its total
energy E is given by
mo c 2
E = mc 2 =
1 - v 2 /c 2
Its momentum is given by
mo v
p = mv =
1 - v 2 /c 2
The frequency of the associated deBroglie wave is given by
E mo c 2 1
n= =
h 2 2 h
1 - v /c
2p m0 c 2
w = 2pn =
h 1 - v 2 /c 2
dw 2p m0 c 2 È 1 Ê 2v ˆ ˘
or = Í- 2 ¥ ÁË - 2 ˜¯ ˙
dv h(1 - v 2 /c 2 )3/2 Î c ˚
dw 2p m0 v
or =
dv h(1 - v 2 /c 2 )3/2
The wavelength of the associated deBroglie wave is given by
h h(1 - v 2 /c 2 )1/2
l= =
p m0 v
Hence propagation constant

2p 2p m0 v
k= =
l h(1 - v 2 /c 2 )1/2
1/2 -1/2
Ê v2 ˆ 1Ê v2 ˆ Ê 2v ˆ
Á 1- 2˜ - v Á1 - ˜ ÁË - 2 ˜¯
dk 2p m0 Ë c ¯ 2Ë c2 ¯ c
or =
dv h Ê v ˆ
2
ÁË1 - 2 ˜¯
c
dk 2p m0 ÈÊ 2 ˆ -1/2 v2 Ê
-3/2 ˘
or = ÍÁ1 - v ˜ v2 ˆ ˙
+ 2 Á1 - ˜
dv h ÎÍË c2 ¯ c Ë c2 ¯ ˚˙
-3/2
2p m0 Ê v2 ˆ È v2 v2 ˘
= Á 1- 2˜ Í1 - 2 + 2 ˙
h Ë c ¯ Î c c ˚
-3/2
2p m0 Ê v2 ˆ
= 1 -
h ÁË ˜
c2 ¯
dk 2p m0
=
dv h(1 - v 2 /c 2 )3/2
Since the group velocity

d w d w /dv
G= =
dk dk /dv
2p m0 v
h(1 - v 2 /c 2 )3/2
G= =v
2p m0
h(1 - v 2 /c 2 )3/2
G = v = the particle velocity
Thus the wave group associated with the moving material particle travels with the same velocity as the
particle. It proves that a material particle in motion is equivalent to group of waves or a wave packet.
s UmmarY

✦ Insufficiency of classical mechanics in explaining phenomena like photoelectric effect, Compton
effect, Raman effect, etc. led to the development of quantum mechanics.
✦ Quantum mechanics developed in two stages. The first stage began with Max Planck’s hypothesis
according to which the radiation is emitted or absorbed by matter in discrete packets or quanta of
energy. This hypothesis led to a theory which was not completely satisfactory being a mixture of
classical and nonclassical concepts. The second stage led to two types of mechanics, namely matrix
mechanics and wave mechanics.
✦ Blackbody radiation shows that i) the distribution of frequencies is a function of temperature of the
blackbody, ii) with the increase in temperature, the total amount of emitted radiation increases, and
iii) the position of the maximum peak shifts toward higher frequencies with increasing equilibrium
temperature.
✦ Planck explained blackbody radiation by using the Maxwell-Boltzmann distribution, i.e., Nn = N0e–E/kT.
He obtained the following expression (called Planck’s radiation law) for the distribution of the maximum
intensity of radiation in the spectrum of blackbody
8p hn 3 1
I (n ) =
c3 e( kv / kT ) - 1
✦ As per Planck’s quantum hypothesis, the atoms of the wall of a blackbody behave as oscillators and
each has a characteristic frequency of oscillation. Then average energy of these Planck’s oscillators
were calculated and finally Planck’s radiation formula was derived.
✦ Wien’s law is deduced from Planck’s radiation formula under the condition when the wavelength l
and temperature T are very small. However, Planck’s radiation formula under the condition of high
temperature T and large wavelength l takes the form of Rayleigh-Jeans law.
✦ Photon is an elementary particle that is massless and has no charge. It is a bundle of energy or packet
of energy emitted by a source of radiation. It moves with velocity of light. It can carry energy and
momentum. If a photon of frequency n is to be treated as a particle, then the characteristics of the
photon are given as
m0 = 0, E = hn, m = hn/c2 and p = hn/c.
✦ Light or electromagnetic radiation exhibits wave and particle properties both, it means it has dual
nature, i.e., it behaves like a particle as well as a wave. The dual characteristic property of radiation is
called wave particle duality.
✦ deBroglie suggested that the dual nature of electromagnetic radiation may be extended to material
particles such as electrons, protons, neutrons etc. It means that a moving particle, whatever its nature be,
has wave properties associated with it. The waves associated with these particles are known as matter
waves or deBroglie waves. The Davisson-Germer experiment on electron diffraction demonstrated
these waves.
h
✦ If there is a particle of momentum p, its motion is associated with a wave of wavelength l = .
p
✦ If there is a wave of wavelength l, the square of the amplitude of the wave at any point in space is
h
proportional to the probability of observing, at that point in space, a particle of momentum p = .
l
✦ The photoelectric effect refers to the emission or ejection of electrons from the surface of a metal
(generally) in response to incident light. Energy contained within the incident light is absorbed by
the electrons within the metal, gaining sufficient energy to be ‘knocked’ out of, i.e., emitted from, the
surface of the metal. The photoelectric effect is perhaps the most direct and convincing evidence of
the existence of photons and the ‘corpuscular’ nature of light and electromagnetic radiation. That is, it
provides undeniable evidence of the quantisation of the electromagnetic field and the limitations of the
classical field equations of Maxwell.
✦ When an electromagnetic radiation (frequency n) is incident on free charges (say, electrons), the free
charges absorb this radiation and start oscillating at frequency n. Then these oscillating charges radiate
electromagnetic waves of the same frequency n. This type of scattering where the change in frequency
or wavelength does not take place is called coherent scattering. This coherent scattering has been
observed with the radiation in visible range and also at longer wavelengths.
✦ In the case of scattering of radiation of very short wavelengths like X-rays, the scattered rays are
found to consist of two frequencies: n and n1. The wavelength l corresponding to the frequency n is
called unmodified wavelength or unmodified radiation, whereas the wavelength l1 corresponding to
the frequency n1 is called modified wavelength or modified radiation. This type of scattering is known
as incoherent scattering.
✦ Compton effect is not significantly observable with visible light, as the Compton shift is extremely
small.
s olved e XamPles
E xamplE 1 Calculate the frequency and wavelength of a photon whose energy is 75 eV.
Solution Given energy E = 75 eV = 75 ¥ 1.6 ¥ 10–19 J.
Formula used is
hc
E = hn =
l
E 75 ¥ 1.6 ¥ 10-19
Frequency (v) = =
h 6.62 ¥ 10-34
= 18.13 ¥ 1015 Hz
c 3 ¥ 108
and wavelength l= =
n 18.13 ¥ 1015
= 1.655 ¥ 10-8 m
= 165.5 ¥ 10-10 m
or l = 165.5 Å
E xamplE 2 Find the number of quanta of energy emitted per second if a radio station operates at a frequency
of 98 MHz and radiates power of 2 ¥ 105 W.
Solution Given n = 98 ¥ 106 cycles/sec and Power (P) = 2 ¥ 105 W = 2 ¥ 105 J/sec.
Energy of each quanta is
E = hn
\ E = 6.62 ¥ 10–34 ¥ 98 ¥ 106
= 6.4876 ¥ 10–26 J/quanta
= 6.5 ¥ 10–26 J/quanta
Number of quanta emitted per second
Power
=
quantum energy
2 ¥ 105 (J/ sec)
6.5 ¥ 10-26 (J/quanta)
= 3.08 ¥ 1030 quanta/sec
E xamplE 3 A certain spectral line has wavelength 4000 Å. Calculate the energy of the photon.
Solution Given l = 4.0 ¥ 10–7 m.
Formula used is
hc
Ek = hn =
l
6.62 ¥ 10-34 ¥ 3 ¥ 108
=
4 ¥ 10-7 m
19
4.965 10 J
E xamplE 4 Calculate the number of photons of green light of wavelength 5000 Å require to make one erg
of energy.
Solution Given l = 5 ¥ 10–7 m.
Formula used is
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108
E= =
l 5 ¥ 10-7
= 3.972 ¥ 10–19 J
= 3.972 ¥ 10–12 erg
Number of photons of green light emitted (per energy)
1.0
=
3.972 ¥ 10-12
252 109
E xamplE 5 Calculate the wavelength of a photon of energy 5 ¥ 10–19 J.

Solution Given E = 5 ¥ 10–19 J
Formula used is
hc
E=
l
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108

or l= =
E 5 ¥ 10-19
= 3.972 ¥ 10-7 m
= 4000 Å
E xamplE 6 Calculate the energy of an electron of wavelength 4.35 ¥ 10–7 m.

Solution Given wavelength (l) = 4.35 ¥ 10–7 m.
Formula used is
hc
E = hv = .
l
6.62 ¥ 10-34 ¥ 3 ¥ 108
=
4.35 ¥ 10-7
19
= 4.566 10 J
E xamplE 7 How many watts of power at the threshold is received by the eye, if it receives 120 photons per
second of the visible light of wavelength = 5600 Å.
Solution Given l = 5.6 ¥ 10–7 m and number of photons = 120.
hc
Energy of a photon E = hn =
l
6.62 ¥ 10-34 ¥ 3 ¥ 108
or E= = 3.55 ¥ 10-19 J
5.6 ¥ 10-7
The energy received by the eye per second = 3.5464 ¥ 120 J/sec
Because 3.55 ¥ 120 = 426
3.5464 ¥ 120 J/sec = 425.57 W = 4.2557 ¥ 10–17 W
E xamplE 8 How many photons of yellow light of wavelength 5500 Å constitute 1.5 J of energy.
Solution Given l = 5.5 ¥ 10–7 m and energy of n photons = 1.5 J
hc
Formula used is E = hn =
l
Energy of a photon of yellow light, i.e.,
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108

E= =
l 5.6 ¥ 10-7
= 3.55 ¥ 10-19 J
Given
n ¥ energy of one photon = 1.5 J
1.5
or n=
3.55 ¥ 10-19
= 4.2 1018
E xamplE 9 Calculate the work function, stopping potential and maximum velocity of photoelectrons for
a light of wavelength 4350 Å when it incidents on sodium surface. Consider the threshold wavelength of
photoelectrons to be 5420 Å.
Solution Given l0 = 5.42 ¥ 10–7 m and l = 4.35 ¥ 10–7 m.
Formulae used are
hc
f0 = = hv0
l0
1 2 È1 1 ˘
mvmax = hc Í - ˙ and
2 Î l l0 ˚
È1 1 ˘
eV = hn - hn 0 = hc Í - ˙
Î l l0 ˚
1 2
or eV = mvmax = (Ek ) max
2
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108
f0 = = = 3.664 ¥ 10-19 J
l0 5.42 ¥ 10-7
1 2 È1 1˘
mvmax = hc Í - ˙
2 Î l l 0˚
2 2hc È l0 - l ˘
vmax =
m ÍÎ l0 l ˙˚
2 ¥ 6.62 ¥ 10-34 ¥ 3 ¥ 108 È (5.42 - 4.35) ¥ 10-7 ˘

= Í -14 ˙
= 0.1981 ¥ 1012
9.1 ¥ 10-31 Î 5.42 ¥ 4.35 ¥ 10 ˚
\ vmax = 0.445 ¥ 106 m/sec
= 4.45 ¥ 105 m/sec
1 2
eV = mvmax
2
The stopping potential
2
mvmax 9.1 ¥ 10-31 ¥ (4.45 ¥ 105 ) 2
V= =
2e 2 ¥ 1.6 ¥ 10-19
0.56 volts
E xamplE 10 The threshold frequency for photoelectric emission in copper is 1.1 ¥ 1015 Hz. Find the maximum
energy in eV when light of frequency 1.2 ¥ 1015 Hz is directed on the copper surface.
Solution Given n0 = 1.1 ¥ 1015 Hz and n = 1.2 ¥ 1015 Hz.
Formula used is
1 2
mvmax = hn - hn 0 = h(n - n 0 )
2
= 6.62 ¥ 10-34 [1.2 - 1.1] ¥ 1015
= 0.662 ¥ 10-19 J
0.414eV
E xamplE 11 Calculate the work function in electron volts of a metal, given that photoelectric threshold is
(i) 6200 Å (ii) 5000 Å.
Solution Given (i) l0 = 6.2 ¥ 10–7 m (ii) l0 = 5.0 ¥ 10–7 m.
hc
f0 = hn 0 =
l0
6.62 ¥ 10-34 ¥ 3 ¥ 108

(i) f0 = eV
6.2 ¥ 10-7 ¥ 1.6 ¥ 10-19
2.0 eV
(ii) l0 = 5.0 ¥ 10-7 m
6.62 ¥ 10-34 ¥ 3 ¥ 108
f0 =
5 ¥ 10-7 ¥ 1.6 ¥ 10-19
= 2.483 eV
2.48 eV
E xamplE 12 Find out the maximum energy of the photoelectron, work function and threshold frequency
when a light of wavelength 3132 Å is incident on a surface of cesium and the stopping potential for the photo
electron is 1.98 volt.
Solution Given V = 1.98 volts and l = 3.132 ¥ 10–7 m.
Formulae used are
1 2
Ek = mvmax = eV0 , V0 = stopping potential
2
Ê1 1ˆ
and Ek = h(n – n 0 ) = hc Á - ˜
Ë l l0 ¯
Then maximum energy of the photoelectron (Emax)
= eV0 = 1.6 ¥ 10-19 ¥ 1.98 J

Ek = 3.168 ¥ 10-19 J
Ê1 1ˆ
Ek = hc Á - ˜
Ë l l 0¯
È 1 1˘
3.168 ¥ 10-19 = 6.62 ¥ 10-34 ¥ 3 ¥ 108 Í -7
- ˙
Î 3.132 ¥ 10 l 0˚
1 1 1
= -
6.2689 ¥ 10-7 3.132 ¥ 10-7 l0
1
or = 3.193 ¥ 106 - 1.595 ¥ 106 = 1.598 ¥ 106
l0
1
l0 = = 6258 Å
1.598 ¥ 106
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108
Work function (f0 ) = =
l0 6.258 ¥ 10-7
= 3.174 ¥ 10-19 J
E xamplE 13 Is it possible to liberate an electron from a metal surface having work function 4.8 eV with an
incident radiation of wavelength (i) 5000 Å and (ii) 2000 Å.
Solution Given f0 = 4.8 eV.
hc
Formula used is Ek = .
l
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108
(i) Energy ( Ek ) = = J
l 5 ¥ 10-7
6.62 ¥ 10-34 ¥ 3 ¥ 108
= eV
5 ¥ 10-7 ¥ 1.6 ¥ 10-19
= 2.48 eV
From the above it is clear that the energy corresponding to wavelength 5000 Å is found to be less than the work function
i.e., 4.8 eV. So it will not be able to liberate an electron.
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108
(ii) Ek = = = 9.93 ¥ 10-19 J
l 2.0 ¥ 10-7
9.93 ¥ 10-19
or Ek = eV = 6.206 eV
1.6 ¥ 10-19
Ek = 6.21 eV
As the energy corresponding to wavelength 2000 Å is greater than the work function. So it is sufficient to liberate
electrons.
E xamplE 14 Find the maximum energy of the photoelectron, the work function and threshold frequency, if
the potassium surface is illuminated by a light of wavelength 5893 Å. The stopping potential for the emitted
electron is 0.36 V.
Solution Given stopping potential V0 = 0.36 V and l = 5893 Å.
Formula used is
Ek = eV = hn – f0
Ek = eV = 0.36 eV
Work function
hc
(f0 ) = hv - eV = - eV
l
6.62 ¥ 10-34 ¥ 3 ¥ 108
= - 0.36 eV
5.893 ¥ 10-7 ¥ 1.6 ¥ 10-19
= 2.11 - 0.36 = 1.75 eV
Thus the work function is 1.75 eV.
Threshold frequency
f
v0 =
h
1.75 ¥ 1.6 ¥ 10-19
=
6.62 ¥ 10-34
= 4.23 ¥ 1014 cycles/sec
E xamplE 15 Find the maximum kinetic energy of the emitted electrons and the stopping potential if the light
of wavelength 5890 Å is incident on the surface for which threshold frequency is 7320 Å.
Solution Given l = 5.89 ¥ 10–7 m and l0 = 7.32 ¥ 10–7 m.
Formula used is
Ê1 1ˆ
Ek = hv - hv0 = hc Á - ˜
Ë l l0 ¯
6.62 ¥ 10-34 ¥ 3 ¥ 108 È 1 1 ˘
= ÍÎ 5.89 - 7.32 ˙˚
10-7
È 7.32 - 5.89 ˘
= 19.86 ¥ 10-19 Í ˙
Î 5.89 ¥ 7.32 ˚
= 6.587 ¥ 10-20 J
Ek
eV = Ek or V =
e
6.587 ¥ 10-20
Stopping potential (V ) =
1.6 ¥ 10-19
= 0.412 V
E xamplE 16 The threshold wavelength for photoelectric emission in Tungsten is 2300 Å. What wavelength
of light must be used in order for electrons with a maximum energy of 1.5 eV to be ejected?
Solution Given l0 = 2.3 ¥ 10–7 m and Ek = 1.5 eV.
Formula used is
Ê1 1ˆ
E = hv - hv0 = hc Á - ˜
Ë l l 0¯
1 1 E
- =
l l0 hc
1 1.5 ¥ 1.6 ¥ 10-19 1
or = +
l 6.62 ¥ 10-34 ¥ 3 ¥ 108 2.3 ¥ 10-7
= 1.2085 ¥ 106 + 4.3478 ¥ 106
1
= 5.556 ¥ 106
l
or l = 1.7998 ¥ 10-7 m
l = 1799.8 Å
E xamplE 17 The work function of Tungsten is 4.53 eV. If ultraviolet light of wavelength 1500 Å is incident
on the surface, does it cause photoelectron emission? If so, what is the kinetic energy of the emitted electron?
Solution Given work function f0 = 4.53 eV and l = 1.5 ¥ 10–7 m.
hc
Formula used is Ek =
l
Energy corresponding to incident photon of wavelength 1.5 ¥ 10–7 m
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108

Ek = = J
l 1.5 ¥ 10-7
= 13.24 ¥ 10-19 J
Ek = 8.28 eV
The kinetic energy of the electron

1 2
Ek = mvmax = hv - hv0 = hv - f0
2
= 8.28 - 4.53
= 3.75 eV
E xamplE 18 The work function of sodium metal is 2.3 eV. What is the longest wavelength of light that cause
photoelectric emission from sodium?
Solution Given f0 = 2.3 eV = 2.3 ¥ 1.6 ¥ 10–19 J.
hc
f0 = hv0 =
l0
Longest wavelength = Threshold wavelength

hc 6.62 ¥ 10-34 ¥ 3 ¥ 108
l0 = =
f0 2.3 ¥ 1.6 ¥ 10-19
= 5396.74 Å
E xamplE 19 Evaluate the threshold wavelength of photoelectric material whose work function is 2.0 eV.
Solution Given f0 = 2.0 eV = 2 ¥ 1.6 ¥ 10–19 J.
Formula used is
hc
l=
f0
6.62 ¥ 10-34 ¥ 3 ¥ 108
or l =
2.0 ¥ 1.6 ¥ 10-19
= 6206 Å
E xamplE 20 Calculate the threshold wavelength and the wavelength of incident electromagnetic radiation
so that the photoelectrons emitted from potassium have a maximum kinetic energy of 4eV. Take the work
function of potassium as 2.2 eV.
Solution Given Emax = 4.0 ¥ 1.6 ¥ 10–19 J and f0 = 2.2 ¥ 1.6 ¥ 10–19 J.
Formulae used are
hc
f0 = hv0 = and Ek = hv - f0 = hv - hv0
l0
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108

l0 = =
f0 2.2 ¥ 1.6 ¥ 10-19
l0 = 5642 Å (Threshold wavelength)
È1 1 ˘
Ek = 4 ¥ 1.6 ¥ 10-19 = hc Í - ˙
Î l l0 ˚
1 4 ¥ 1.6 ¥ 10-19 1
or = +
l 6.62 ¥ 10-34 ¥ 3 ¥ 108 l0
= 3.223 ¥ 106 + 1.772 ¥ 106
1
= 4.995 ¥ 106
l
1
l=
4.995 ¥ 106
= 2002 Å
E xamplE 21 Ultraviolet light of wavelength 350 nm and intensity 1.0 watt/m2 is directed at a potassium
surface. (i) Find the maximum kinetic energy of photoelectron (ii) 0.5% of incident photons produce
photoelectrons, how many photoelectrons are emitted per second if the surface of potassium is 1.0 cm2. Work
function of potassium is 2.1 eV.
Solution Given l = 3.5 ¥ 10–7 m and f0 = 2.1 eV.
(i) Formula used is
1 2
Ek = mvmax = hv - f0 .
2
hc 6.62 ¥ 10-34 ¥ 3 ¥ 108
= - f0 = - 2.1 eV
l 3.5 ¥ 10-7 ¥ 1.6 ¥ 10-19
Ek = (3.546 - 2.1) eV = 1.45 eV
= 2.3136 ¥ 10-19 J
= 2.314 ¥ 10-19 J
(ii) Energy incident per second on 1.0 cm2 surface of potassium = 10–4 Joule
The energy which produces photoelectron per second = 0.5%.
0.5
Effective energy which will be used to produce photoelectrons = ¥ 10-4 J = 5 ¥ 10-7 J
100
Minimum energy required to eject one electron from the surface
= 2.314 ¥ 10–19 J
5 ¥ 10-7
So the number of electrons emitted per second from 1.0 cm2 area of the surface of potassium will be = .
2.314 ¥ 10-19
= 2.16 × 1012
E xamplE 22 Calculate the value of Planck’s constant from the following data, assuming that the electronic charge
e has value of 1.6 ¥ 10–19 Coulomb. A surface when irradiated with light of wavelength 5896 Å emits electrons for
which the stopping potential is 0.12 volts. When the same surface is irradiated with light of wavelength 2830 Å, it
emits electrons for which the stopping potential is 2.2 volts.
Solution If the radiation of wavelength is incident on the surface of the metal having work function f0 and stopping
potential V0 for the emitted electrons, then f0 and V0 satisfy the following relation.
hc
eV0 = - Ve f0 (i)
l
(i) Given l = 5.896 ¥ 10–7 m and V0 = 0.12 volts
hc
= eV0 + f0
l
h ¥ 3 ¥ 108
= 1.6 ¥ 10-19 ¥ 0.12 + f0 (ii)
5.896 ¥ 10-7
(ii) Given l = 2.83 ¥ 10–7 m and V0 = 2.2 volts, then
h ¥ 3 ¥ 108
= 1.6 ¥ 10-19 ¥ 2.2 + f0 (iii)
2.83 ¥ 10-7
On subtracting Eq. (ii) from Eq. (iii), we get
È 3 ¥ 108 3 ¥ 108 ˘ -19 -19
hÍ -7
- ˙ = ÈÎ1.6 ¥ 10 ¥ 2.2 - 1.6 ¥ 10 ¥ 0.12˘˚
Î 2.83 ¥ 10 5.896 ¥ 10-7 ˚
3 ¥ 1015[5.896 - 2.83]
h¥ = 1.6 ¥ 10-19 ¥ 2.08
2.83 ¥ 5.896
h = 6.04 ¥ 10-34 Jsec
E xamplE 23 Calculate Compton shift if X-rays of wavelength 1.0 Å are scattered from a carbon block. The
scattered radiation is viewed at 90° to the incident beam.
Solution Given l = 1.0 Å = 10–10m and f = 90º.
Formula used is
h
Dl = (1 - cos f )
m0c
6.62 ¥ 10-34
= (1 - cos 90∞)
9.1 ¥ 10-31 ¥ 3 ¥ 108
= 0.242 ¥ 10-11 m
= 0.024 ¥ 10-10 m
= 0.0242 Å
E xamplE 24 An X-ray photon is found to have doubled its wavelength on being scattered by 900. Find the
energy and wavelength of incident photon.
Solution Given f = 90º.
Formula used is
h
Dl = (1 - cos f )
m0c (i)
6.62 ¥ 10-34
= (1 - cos 90∞)
9.1 ¥ 10-31 ¥ 3 ¥ 108
= 0.242 ¥ 10-11 m = 0.024 Å
As Dl = l¢ – l, where l is the wavelength of incident photon and l¢ is the wavelength of scattered photon, then
l¢ = l + Dl (ii)
Given l¢ = 2l (iii)
From Eqs. (ii) and (iii), we get
2l = l + Dl
or l = Dl = 0.0242 ¥ 10–10 m = 0.0242 Å
hc
Energy of the incident photon (E) = hv =
l
6.62 ¥ 10-34 ¥ 3 ¥ 108
= = 0.513 MeV
0.0242 ¥ 10-10
E xamplE 25 Calculate the wavelength of incident X-ray photon which produces recoil electron of energy
4.0 KeV in Compton effect. The electron recoils in the direction incident photon and photon is scattered at
an angle of 180º.
Solution f = 180º and energy of the recoiled electron = 4000 eV.
Let l be the wavelength of incident X-ray photon and l¢ be the scattered photon, then according to the law of conservation
of energy.
hc hc
- = Kinetic energy of the recoiled electron
l l¢
1 2
= mv = 4 ¥ 103 eV = 4 ¥ 103 ¥ 1.6 ¥ 10-19 J
2
hc hc (i)
or - = 6.4 ¥ 10-16 J
l l¢
According to the principal of conservation of linear momentum in the direction incident photon
h h
= cos f + mv cos q
l l¢ (ii)
h h
= cos 180∞ + mv cos 0∞ = - + mv
l¢ l¢
h h
+ = mv (iii)
l l¢
Momentum (p = mn) can be calculated as
1 2
mv = 4.0k eV = 4 ¥ 103 ¥ 1.6 ¥ 10-19 J
2
= 6.4 ¥ 10-16 J
m
mv 2 = 2m ¥ 6.4 ¥ 10-16 J = 2 ¥ 9.1 ¥ 10-31 ¥ 6.4 ¥ 10-16
m
= 1.1648 ¥ 10-45 = 11.648 ¥ 10-46
mv = 34.13 ¥ 10-24 kgm sec -1 (iv)
By using Eqs. (iii) and (iv) then, we get
h h
+ = 34.13 ¥ 10-24
l l¢
Multiplying by velocity of light
hc hc
+ = 102.4 ¥ 10-16 (v)
l l¢
By adding Eq. (v) with Eq. (i), we get
hc
2 = (102.4 + 6.4) ¥ 10-16
l
= 108.79 ¥ 10-16
2hc
l= = 0.365 ¥ 10-10 m
108.79 ¥ 10-16
l = 0.365 Å
E xamplE 26 X-rays with l = 1 Å are scattered from a carbon block. The scattered radiation is viewed at 900
to the incident beam.
1. What is Compton shift Dl?
2. What kinetic energy is imparted to the recoil electron?
Solution Given l = 1 ¥ 10–10 m.
Formula used for Compton shift is
h
Dl = (1 - cos f )
m0c
f = 90∞
6.62 ¥ 10-34
Dl = (1 - cos 90∞)
9.1 ¥ 10-31 ¥ 3 ¥ 108
= 2.425 ¥ 10-12 m
Let l be the wavelength of incident X-ray photon and l¢ be the scattered photon, then according to the law of conservation
of energy
hc hc hc
= + Ek = + Ek
l l¢ l + Dl
hc hc hcDl
Ek = - =
l l + Dl l (l + Dl )
6.62 ¥ 10-34 ¥ 3 ¥ 108 ¥ 2.425 ¥ 10-12
=
1 ¥ 10-10 ¥ (1 + 0.02425) ¥ 10-10
= 47.02 ¥ 10-18 J
= 294 eV
E xamplE 27 X-ray of wavelength 0.144 Å are scattered from a carbon target. Find maximum shift in
wavelength and maximum energy of recoil electron.

Formula used for Compton shift
h
Dl = (1 - cos f )
m0c
The Compton shift will be maximum if f = 180º
6.62 ¥ 10-34
[ Dl ]max = (1 - cos 180∞)
9.1 ¥ 10-31 ¥ 3 ¥ 108
2 ¥ 6.62 ¥ 10-34
=
9.1 ¥ 10-31 ¥ 3 ¥ 108
= 0.485 ¥ 10-11 m
= 0.0485 Å
The kinetic energy of the recoil electron is given by the relation
hc hc Ê1 1ˆ
Ek = - = hc Á - ˜
l l¢ Ë l l¢¯
l ¢ = l + Dl , then
È1 1 ˘
Ek = hc Í - ˙
Î l l + Dl ˚
6.62 ¥ 10-34 ¥ 3 ¥ 108 È 1 1 ˘
= Í 0.144 - 0.144 + 0.0485 ˙
10-10 Î ˚
= 34.75 ¥ 10-16 J
= 21.72 keV
E xamplE 28 X- rays of wavelength 0.2 Å are scattered from a target. Calculate the wavelength of X-ray
scattered through 45º. Also find the maximum kinetic energy of the recoil electron.
Solution Given l = 0.2 Å = 0.2 ¥ 10–10 m and f = 45º.
h 6.62 ¥ 10-34
Dl = (1 - cos 45∞) = [1 - 0.7071] = 0.0071 Å
mc 9.1 ¥ 10-31 ¥ 3 ¥ 108
Therefore, wavelength of scattered X-rays
l ¢ = l + Dl = 0.2 + 0.0071 = 0.2071 Å

Ê1 1ˆ
Ek = hc Á - ˜
Ë l l¢¯
Thus kinetic energy is maximum if l¢ is maximum. The maximum value of l¢ can be obtained by the relation l¢ = l + Dlm.
Maximum value of Dl is obtained at f = 180º.
h 6.62 ¥ 10-34
Dlm = (1 - cos f ) = (1 - cos 180∞)
m0c 9.1 ¥ 10-31 ¥ 3 ¥ 108
2 ¥ 6.62 ¥ 10-34
= = 0.0485 Å
9.1 ¥ 10-31 ¥ 3 ¥ 108
l ¢ = 0.2 + 0.0485
= 0.2485 Å
Hence, maximum kinetic energy i.e.,

6.62 ¥ 10-34 ¥ 3 ¥ 108 È 1 1 ˘
Ek = ÍÎ 0.2 - 0.2485 ˙˚
10-10
= 19.38 ¥ 10-16 J
E xamplE 29 Calculate the deBroglie wavelength associated with the automobile of mass 2 ¥ 103 kg which
is moving with a speed 96 km/hr.
96 ¥ 103
Solution Given m = 2 ¥ 103 kg, v = m/sec = 26.67 m/sec.
60 ¥ 60
deBroglie wavelength is given as
h 6.62 ¥ 10-34
l= =
mv 2 ¥ 103 ¥ 26.67
= 0.124 ¥ 10-37 m
= 1.24 ¥ 10-38 m
E xamplE 30 A particle of charge q and mass m is accelerated through a potential difference V. Find its deBroglie
wavelength. Calculate the wavelength (l), if the particle is an electron and V = 50 volts.
Solution When a particle of charge q and mass m is accelerated through a potential V, then deBroglie wavelength is
given by h
l=
mv (i)
1
and Ek = mv 2 = qV 2 2
or m v = 2mqV
2
or mv = 2mqV (ii)
By using Eqs. (i) and (ii), we obtain
h
l=
2mqV
Given q = 1.6 ¥ 10–19 C and V = 50 volts, then
6.62 ¥ 10-34
l=
2 ¥ 9.1 ¥ 10-31 ¥ 1.6 ¥ 10-19 ¥ 50
= 1.74 Å
E xamplE 31 Calculate the wavelength of thermal neutrons at 27ºC, given mass of neutron = 1.67 ¥ 10–27 kg,
Planck’s constant h = 6.6 ¥ 10–34 J sec and Boltzmann’s constant k = 1.37 ¥ 10–23 JK–1.
Solution Given T = 27ºC = 27 + 273 = 300K, m = 1.67 ¥ 10–27kg, h = 6.6 ¥ 10–34 Jsec and k = 1.376 ¥ 10–23JK–1.
deBroglie wavelength is given by
h
l= (i)
mv
1 3
Et = mv 2 = kT or (mv) 2 = 3mkT
2 2
or mv = 3mkT (ii)
h 6.6 ¥ 10-34
Then, l = =
3mkT 3 ¥ 1.67 ¥ 10-27 ¥ 1.376 ¥ 10-23 ¥ 300
= 1.452 ¥ 10-10 m
or l = 1.452 Å
ExamplE 32 A proton is moving with a speed 2 ¥ 108 m/sec. Find the wavelength of matter wave associated
with it.
Solution Given v = 2 ¥ 108 m/sec.
Formula used for deBroglie wavelength is
6.62 ¥ 10-34
l=
1.67 ¥ 10-27 ¥ 2 ¥ 108
= 1.98 ¥ 10-15 m
E xamplE 33 The deBroglie wavelength associated with an electron is 0.1 Å. Find the potential difference
by which the electron is accelerated.
deBroglie wavelength in terms of potential difference is given by
h
l=
2mqV
h2
or 2mqV =
l2
h2 (6.62 ¥ 10-34 ) 2
V= =
2mql 2
2 ¥ 9.1 ¥ 10-31 ¥ 1.6 ¥ 10-19 ¥ (10-11 ) 2
= 15.05 kV
E xamplE 34 Calculate the deBroglie wavelength of an a-particle accelerated through a potential difference
of 200 volts.
Solution Given V = 200 volts, q = qa = 2e = 3.2 ¥ 10–19 C and m = ma = 4mp.
deBroglie wavelength in terms of potential difference
h 6.62 ¥ 10-34
l= =
2ma qV 2 ¥ 4 ¥ 1.67 ¥ 10-27 ¥ 2 ¥ 1.6 ¥ 10-19 ¥ 200
6.62 ¥ 10-34
= = 0.07159 ¥ 10-11
92.468 ¥ 10-23
l = 7.16 ¥ 10-13 m
E xamplE 35 Calculate the deBroglie wavelength of an average Helium atom in furnace of 400 K. Given
k = 1.38 ¥ 10–23 J/K
Solution Given T = 400 K, k = 1.38 ¥ 10–23 J/K and mass of Helium atom = 4mp = 4 ¥ 1.67 ¥ 10–27 kg.
deBroglie wavelength in terms of temperature i.e.,
Engineering Physics Development of Quantum Mechanics 585
h 6.62 ¥ 10-34
l= =
3mkT 3 ¥ 4 ¥ 1.67 ¥ 10-27 ¥ 1.38 ¥ 10-23 ¥ 400
6.62 ¥ 10-34
= = 0.6294 Å
105.176 ¥ 10-25
l = 0.6294 Å
E xamplE 36 Calculate the deBroglie wavelength associated with a neutron moving with a velocity of
2000 m/sec.
Solution Given v = 2000 m/sec and m = 1.67 ¥ 10–27 kg.
deBroglie wavelength
h 6.62 ¥ 10-34
l= =
mv 1.67 ¥ 10-27 ¥ 2000
= 1.98 ¥ 10-10 m
= 1.98 Å
E xamplE 37 Calculate the energy in eV corresponding to a wavelength of 1.0 Å for electron and neutron.
Given h = 6.6 ¥ 10–34 J sec, mass of electron = 9.1 ¥ 10–31 kg and mass of the neutron = 1.7 ¥ 10–27 kg.
Solution Formula used is
h h
l= or
mv l m
6.6 ¥ 10-34
or v=
1.0 ¥ 10-10 ¥ 1.7 ¥ 10-27
= 3.88 ¥ 103 m/sec
If the velocity is much less than the velocity of light, it can be considered as non-relativistic case and hence deBroglie
wavelength can be obtained by the relation.
h h2
l= or l 2 =
2mE 2mE
For Electron
h2 (6.62 ¥ 10-34 ) 2
E= =
2ml 2
2 ¥ 9.1 ¥ 10-31 ¥ (10-10 )2
43.8244 ¥ 10-68
= = 2.41 ¥ 10-17 J
18.2 ¥ 10-51
=1.51 ¥ 100 = 151 eV
E = 151 eV
For neutron
(6.62 ¥ 10-34 ) 2 43.8244 ¥ 10-68

E= -27 -10 2
=
2 ¥ 1.7 ¥ 10 ¥ (10 ) 3.4 ¥ 10-47
= 12.89 ¥ 10-21 J
= 0.081 eV
E xamplE 38 Calculate deBroglie wavelength of an electron whose kinetic energy is (i) 500 eV, (ii) 50 eV
and (iii) 1.0 eV.
h
l=
2mE
(i) E = 500 eV = 500 ¥ 1.6 ¥ 10–19 = 8.0 ¥ 10–17 J
6.62 ¥ 10-34
l= = 5.486 ¥ 10-11 m
-31 -17
2 ¥ 9.1 ¥ 10 ¥ 8.0 ¥ 10
= 0.5486 Å
(ii) E = 500 eV = 50 ¥ 1.6 ¥ 10–19 = 8.0 ¥ 10–18
6.62 ¥ 10-34
l=
2 ¥ 9.1 ¥ 10-31 ¥ 8 ¥ 10-18
l = 1.735 ¥ 10-10 m
or l = 1.735 Å
(iii) E = 1.0 eV = 1.6 ¥ 10–19 J
6.62 ¥ 10-34
l=
2 ¥ 9.1 ¥ 10-31 ¥ 1.6 ¥ 10-19
l = 12.267 Å
E xamplE 39 Calculate the ratio of deBroglie wavelengths associated with the neutrons with kinetic energies
of 1.0 eV and 510 eV.
h
l=
2mE
For E = 1.0 eV = 1.6 ¥ 10–19 J and mn = 1.7 ¥ 10–27 kg
6.62 ¥ 10-34
l1 =
2 ¥ 1.7 ¥ 10-27 ¥ 1.6 ¥ 10-19
= 2.838 ¥ 10-11
l1 = 0.284 Å
For E = 510 eV = 510 ¥ 1.6 ¥ 10-19 = 816 ¥ 10-19 J
6.62 ¥ 10-34
l2 =
2 ¥ 1.7 ¥ 10-27 ¥ 816 ¥ 10-19
l2 = 0.01257 Å
= 0.0126 Å
and ratio of deBroglie wavelength is
l1 0.284
= = 22.54 : 1
l2 0.0126
E xamplE 40 Calculate the ratio of deBroglie waves associated with a proton and an electron each having the
kinetic energy as 20 M eV [mp = 1.67 ¥ 10–27 kg and me = 9.1 ¥ 10–31 kg].
Solution Given energy of each proton and electron is 20 ¥ 106 ¥ 10–19 J = 3.2 ¥ 10–12 J.
Formula used is
h
l=
3mE
For proton
6.62 ¥ 10-34
lp =
2 ¥ 1.67 ¥ 10-27 ¥ 3.2 ¥ 10-12
= 6.4 ¥ 10-15 m
For electron
6.62 ¥ 10-34
le =
2 ¥ 9.1 ¥ 10-31 ¥ 3.2 ¥ 10-12
= 2.74 ¥ 10-13 m
The ratio of lp to le is
lp : le = 1:43
E xamplE 41 Calculate the deBroglie wavelength of 1.0 M eV proton. Do we require relativistic calculation?
Solution Given Energy E = 1.0 ¥ 106 ¥ 1.6 ¥ 10–19 J = 1.6 ¥ 10–13 J
Formula used for velocity of Proton
1 2 2E
E= mv or v 2 =
2 m
2E 2 ¥ 1.6 ¥ 10-13
or v= =
m 1.67 ¥ 10-27
= 1.38 ¥ 107 m/sec
From the above result it is clear that the velocity of proton is nearly one twentieth of the velocity of light. So the
relativistic calculations are not required.
E xamplE 42 Calculate the deBroglie wavelength associated with a proton moving with a velocity equal to
1/20th of velocity of light.
c 3 ¥ 108
Solution Given v = = = 1.5 ¥ 107 m/sec and m = 1.67 ¥ 10-27 kg
20 20
Formula used is
h h 6.62 ¥ 10-34
l= = =
p mv 1.67 ¥ 10-27 ¥ 1.5 ¥ 107
= 2.643 ¥ 10-14 m
E xamplE 43 Calculate the kinetic energy of a proton and an electron so that the deBroglie wavelengths
associated with them is the same and equal to 5000 Å.
Solution Given wavelength of proton and electron = 5.0 ¥ 10–7 m.

Formula used in
h h2
l= or E=
2mE 2ml 2
For proton m = mp = 1.67 ¥ 10–27 kg and l = 5.0 ¥ 10–7 m
(6.62 ¥ 10-34 ) 2
E=
2 ¥ 1.67 ¥ 10-27 ¥ (5.0 ¥ 10-7 ) 2
43.8244 ¥ 10-68
= = 0.5248 ¥ 10-27 J
83.5 ¥ 10-41
= 5.248 ¥ 10-28 J
For electron m = me = 9.1 ¥ 10–31

(6.62 ¥ 10-34 ) 2 43.8244 ¥ 10-68
E= -31 -7 2
=
2 ¥ 9.1 ¥ 10 ¥ (5 ¥ 10 ) 4.55 ¥ 10-43
E = 9.63 ¥ 10-25 J
E xamplE 44 Find deBroglie wavelength of an electron in the first Bohr’s orbit of hydrogen atom.
-13.6
Solution Energy of an electron in the first Bohr’s orbit of hydrogen atom can be obtained by using the relation En = 2
n
-13.6
E1 = = - 13.6 eV
12
E1 = - 13.6 ¥ 1.6 ¥ 10-19 J = - 2.176 ¥ 10-18 J
Magnitude of energy = 2.176 ¥ 10-18 J
h 6.62 ¥ 10-34
Wavelength l = =
2mE 2 ¥ 9.1 ¥ 10-31 ¥ 21.76 ¥ 10-19
= 3.3 ¥ 10-10 m
= 3.3 Å
E xamplE 45 Calculate the ratio of deBroglie wavelengths of a hydrogen atom and helium atom at room
temperature, when they move with thermal velocities. Given mass of hydrogen atom mH = 1.67 ¥ 10–27 kg and
mass of helium atom mHe = 4 ¥ mp = 4 ¥ 1.67 ¥ 10–27 kg at room temperature T = 27ºC = 300 K and Boltzmann’s
constant k = 1.376 ¥ 10–23 J/K.
Solution deBroglie wavelength can be calculated by the relation
h
l=
3mkT
For Hydrogen atom
6.62 ¥ 10-34
l=
3 ¥ 1.67 ¥ 10-27 ¥ 1.376 ¥ 10-23 ¥ 300
= 1.456 ¥ 10-10 m
l = 1.456 Å
For Helium atom

6.62 ¥ 10-34
lHe =
3 ¥ 4 ¥ 1.67 ¥ 10-27 ¥ 1.376 ¥ 10-23 ¥ 300
= 0.728 ¥ 10-10 m
l = 0.728 Å
The ratio of wavelengths i.e.,
lH 1.456 2
= =
lHe 0.728 1
H: He = 2:1
E xamplE 46 A proton and a deuteron have the same kinetic energy. Which has a longer wavelength?
Solution mp = mass of proton, md = 2mp and vp and vd are the velocities of proton and deuteron.
Kinetic energy of proton is given by
1
Ep = m p v 2p
2
and kinetic energy of deuteron is
1 1
Ed = md vd2 = (2m p )vd2
2 2
2
Ed = m p vd
But E p = Ed , then
1
m p vd2 = m p v 2p
2
vp
or vd =
2
deBroglie wavelength corresponding to moving proton and deuteron are
h
lp = and
m pv p
h h h
ld = = =
md vd 2m p v p 2m p v p
2
ld h m pv p 1
= ¥ =
lp 2m p v p h 2
lp = 2 d
i.e., proton has a longer wavelength.

E xamplE 47 Find the phase and group velocities of an electron whose deBroglie wavelength is 1.2 Å.
h
l= (i)
mv
vg= Group velocity = Particle velocity = v

h 6.62 ¥ 10-34
v= =
ml 9.1 ¥ 10-31 ¥ 1.2 ¥ 10-10
v = 6.06 ¥ 106 m/sec = group velocity = vg
or vg = 6.06 ¥ 106 m/sec
w
Phase velocity v p = (ii)
k
Energy E = hv
h (iii)
or E= 2p v = w
2p
h
and momentum p =
l
h 2p
or p= = k (iv)
2p l
w w E
or vp = = =
k k p
1 2
and E = mv and p = mv (v)
2
m 2v 2 p 2 (vi)
or E= =
2m 2m
E From Eq. (v)
vp =
p
p 2 /2m p h/l h
vp = = = =
p 2m 2m 2ml
6.62 ¥ 10-34
=
2 ¥ 9.1 ¥ 10-31 ¥ 1.2 ¥ 10-10
v p = 3.03 × 106 m/sec
From the above result it is clear that the phase velocity is just half of group velocity.
E xamplE 48 Calculate the deBroglie wavelength of
(a) a particle accelerated by a potential difference of 30,000 V and
(b) an electron moving with a velocity of 0.01c, where c is the speed of light.
Solution Given V = 30,000 Volts, e = 1.6 ¥ 10–19 Coulomb, me = 9.1 ¥ 10–31 kg, h = 6.63 ¥ 10–34 J.sec
and E = eV = 1.6 ¥ 10–19 ¥ 30,000
= 4.8 ¥ 10–15 Joules
(a) Formula used is
h h È 1 2˘
l= = ÍÎ E = 2 mv ˙˚
mv 2mE
6.63 ¥ 10-34
l= = 7.09 ¥ 10-12 m
2 ¥ 9.1 ¥ 10-31 ¥ 4.8 ¥ 10-15
(b) Formula for deBrogile wavelength is given by

h 6.63 ¥ 10-34 6.63 ¥ 10-34
l= = =
mv 9.1 ¥ 10-31 ¥ 0.01 c 9.1 ¥ 10-31 ¥ 0.01 ¥ 3 ¥ 108
= 2.43 Å
E xamplE 49 Calculate the deBreoglie wavelength of virus particle of mass 1 ¥ 10–15 kg moving at a speed
–3
of 2 ¥ 10 m/sec.
Solution Given, v = 2 ¥ 10–3 m/sec; m = 1 ¥ 10–15 kg
h 6.63 ¥ 10-34
l= = = 3.315 ¥ 10-16 m
mv 1 ¥ 10-15 ¥ 2 ¥ 10-3
Q.1 Which of the following phenomena cannot be explained by the classical theory?
(a) Photoelectric effect (b) Compton effect
(c) Raman effect (d) All of these
Q.2 Which of the following phenomena show the particle nature of light?
(a) Photoelectric effect (b) Raman effect
(c) Compton effect (d) All of these
Q.3 Wien’s law is deduced from Planck’s radiation formula under the condition of
(a) very small wavelength and temperature (b) large wavelength and high temperature
(c) small wavelength and high temperature (d) large wavelength and low temperature
Q.4 Rayleigh–Jeans law is deduced from the Planck’s radiation formula under the condition of
(a) large wavelength and high temperature (b) small wavelength and low temperature
(c) small wavelength and high temperature (d) large wavelength and low temperature
Q.5 Which of the following characteristic(s) photon has(have)?
hv
(a) m0 = 0 (b) E = hn (c) m = 2 (d) All of these
c
Q.6 Which of the following relation can be used to determine deBroglie wavelength associated with a
particle of mass m and having energy E
h h h
(a) (b) (c) (d) All of these
2mqV 3mkT 2mE
Q.7 The phase velocity of deBroglie wave associated with an electron is given by
E hc
(a) (b) hv (c) (d) k
p l
Q.8 Electron behaves like a wave as it
(a) can be deflected by an electric field (b) can be deflected by a magnetic field
(c) they ionise a gas (d) can be diffracted by a crystal
Q.9 A material particle is in thermal equilibrium at temperature T. The wavelength of deBroglie wave
associated with it is
h h h h
(a) (b) (c) (d)
2kT 8p 2 mkT 2mkT 4p 2 mkT
Q.10 Photoelectric effect involves only
(a) free-electron (b) bound electron
Q.11 The deBroglie hypothesis is concerned with
(a) wave nature of radiations (b) wave nature of all material particles
(c) wave nature of electrons only (d) wave nature of a-particles only
Q.12 A proton and a deuteron have the same kinetic energy. The relation between the wavelengths of waves
associated with them is
(a) lp > ld (b) ld = lp (c) lp = ld (d) none of these
Q.13 The group velocity of matter waves is
(a) equal to the particle velocity (b) greater than the particle velocity
(c) less than the particle velocity (d) same as phase velocity
Q.14 The ratio of deBroglie wavelengths of a hydrogen atom and helium atom at room temperature, when
they move with thermal velocities, is
(a) 1:2 (b) 2:1 (c) 3:1 (d) 1:3
Q.15 The existence of matter wave is experimentally proved by
(a) Raman (b) Davisson and Germer
(c) deBroglie (d) none of these
Q.16 Dual character of matter was proposed by
(a) Davisson and Germer (b) deBroglie
(c) Planck (d) none of these
Q.17 Quantum theory successfully explains the phenomena of
(a) photoelectric and compton effects (b) interference, diffraction and polarisation
(c) black body radiations (d) all of these
Q.18 Matter waves
(a) show diffraction (b) show interference
(c) polarisation (d) none of these
Q.19 Matter waves are similar in nature to
(a) cathode rays (b) electromagnetic waves
(c) X-rays (d) both (a) & (b)
Q.20 Tick the correct target material in Davisson Germer experiment
(a) coper (b) nickel (c) silver (d) none of these
Q.1 What do you understand by blackbody radiation?

Q.2 What is Planck’s quantum hypothesis?
Q.3 What do you mean by duality of matter?
Q.4 What is photoelectric effect?
Q.5 Explain Compton effect.
Q.6 What do you understand by wave velocity and group velocity in the context of deBroglie waves?
Q.7 What is the difference between phase velocity and group velocity in the context of deBroglie waves?
Q.8 Discuss the shortcomings of the classical physics and also explain how did quantum mechanics
develop?
Q.9 What is Planck’s constant? Discuss its importance.
Q.10 What are the limitations of old quantum theory?
P ractice P roblems
General Questions
Q.1 What are the shortcomings of old quantum theory?
Q.2 Discuss the failures of classical physics and how does quantum mechanics overcome these failures?
Q.3 What is Planck’s quantum hypothesis to explain the observed spectrum of a black body?
Q.4 Explain briefly quantum theory of radiation. What is a photon? State its properties. Express the linear
Æ
momentum of a photon in terms of wave vector |k | and energy of photon in terms of angular velocity w.
Q.5 Discuss and derive Planck’s radiation formula. Explain Wien’s law and Rayleigh–Jeans law as the
special cases of it.
Q.6 What is photoelectric effect? Draw a labelled diagram of the apparatus you will use to demonstrate
photoelectric effect. Write down its important results and show them graphically.
Q.7 (a) State the laws of photoelectric emission.
(b) In what way classical electromagnetic theory of light fails to explain the basic facts of photoelec-
tricity.
Q.8 Derive Einstein’s photoelectric equation. How does it explain the laws of photoelectric emission? Why
all the photoelectrons do not have the same energy?
Q.9 What is meant by work function of a material? Show how you will measure it experimentally?
Q.10 Draw a curve showing stopping potential against frequency of a photosensitive material. How do you
determine the following with the help of the curve. (a) Threshold frequency (b) Work function and (c)
Planck’s constant.
Q.11 Explain the concept of wave particle dualism. What led deBroglie to suggest that matter has wave
characteristic?
Q.12 State the deBroglie hypothesis of matter waves. Derive an expression for deBroglie wavelength of
matter particle in terms of kinetic energy and temperature.
Q.13 Derive a formula expressing deBroglie wavelength of an electron in terms of potential difference (V)
in volts through which it is accelerated.
Q.14 Why are wave properties of particles normally observed only when we study very small particles?
Q.15 Why can’t we observe deBroglie wavelength with a fast moving cricket ball?
Q.16 Describe with necessary theory the Davisson and Germer experiment for establishing wave nature of
the electron?
Q.17 What is the effect of increasing the electron energy on the scattering angle in a Davisson and Germer
experiment?
Q.18 What is Compton effect? Derive an expression for Compton shift and wavelength of scattered photon.
Explain why Compton shift is not observed with visible light?
Q.19 What is Compton wavelength? Determine its value. Distinguish between Compton shift and Compton
wavelength. What are the factors on which Compton shift depends?
Q.20 What is the difference between phase velocity and group velocity?
Q.21 Prove that the wave group associated with a moving particle travels with the same velocity as that of
the particle?
Q.22 Show that group velocity and wave velocity are the same in a non-dispersive medium?
dv
Q.23 Show that the group velocity G = - l 2 , where the symbols have their usual meanings?
dl
Q.24 Explain group velocity and phase velocity. Derive the expression for group velocity with which a wave
group travels?
Q.25 What is the difference between phase and group velocities. Show that the deBroglie group velocity
associated with the wave packet is equal to the velocity of the particle?
Q.26 Distinguish between phase velocity and group velocity. Show that of a non-relativistic free particle
phase velocity is half of the group velocity?
Q.1 Find out the wavelength associated with photon of energy 10–19 J and also find the energy in eV.
[Ans: 19800 Å & 0.63 eV]
Q.2 It source is operating at a frequency of 108 Hz and radiates a power 104 J/sec, what would be the
number of quanta of energy emitted in one second. [Ans: 1.51 ¥ 1030]
Q.3 A 10 kilowatt transmitter operates at a frequency of 880 kHz. How many photons per second are
emitted? [Ans: 1.716 ¥ 1031]
Q.4 How many photons of red light of wavelength 7800 Å constitute 2.0 J of energy? [Ans: 7.85 ¥1018]
Q.5 Calculate the work function in electron volts of a metal, when photoelectric threshold wavelength is
6800 Å. [Ans: 1.83 eV]
Q.6 The threshold wavelength for photoelectric emission in tungsten is 230 nm. What wavelength of
incident light must be used in order to eject electrons with a maximum velocity of 5 ¥ 105 m/sec
[Ans: 203 nm]
Q.7 Light of wavelength 2000 Å falls on a photosensitive material having work function 4.2 eV. What is
the kinetic energy of the fastest and slowest photo-electron? Also calculate the stopping potential.
[Ans: 2.0 eV, Zero, 2.0 V]
Q.8 When X-rays of energy 0.1 MeV strike a target, they are scattered at an angle of 30º, Compute the
energy of X-rays scattered and the energy of recoiled electron. [Ans: 97.44 keV, 2.56 keV]
Q.9 Compute the deBroglie wavelength of an electron whose kinetic energy is 50 eV. (Given
h = 6.62 ¥ 10–34 Jsec, m = 9.1 ¥ 10–31 kg and eV = 1.6¥10–19J.) [Ans: 1.73 Å]
Q.10 Each of a photon and an electron has an energy of 1 keV. Calculate their corresponding wavelengths.
[Ans: 12.4 Å, 0.39 Å]
–14
Q.11 Find the energy of neutron having deBroglie wavelength 10 m. Given rest mass of neutron as 1.6 ¥
10–27 kg. [Ans: 8.5 meV]
Quantum Mechanics 16
Learning Objectives
LO 1 Understand and learn the Heisenberg LO 3 Know about operators associated with
uncertainty principle and its measurable parameters
applications LO 4 Explain applications of Schrödinger
LO 2 Describe how to obtain time equation
independent/dependent Schrödinger LO 5 Discuss quantum statistics
equation
Introduction
The wave like and particle behaviour of electrons and photons have been discussed in the previous chapter.
However, all the subatomic particles like protons, neutrons, etc. show their dual nature, i.e., sometimes they
behave as particle and sometimes as wave. Various types of explanation to understand this wave particle
duality led to the development of quantum mechanics. Quantum mechanics deals with the behaviour and
characteristics of matter, in the subatomic level, and energy. With the development of quantum theory,
queries like stability of electron orbits and blackbody radiation could be explained scientifically.
Basics of quantum theory were developed by Planck, Einstein, Schrödinger and Heisenberg. As discussed
earlier, Planck in 1900 established that all forms of matter emit or absorb energy in units, called quanta.
Prior to this theory, it was assumed that energy existed only in the form of electromagnetic waves. In
1905, Einstein stated that not only energy but also radiation is quantifiable. He came to the conclusion that
the energy (E) of light depends on its frequency (n) as per the relation E = hn. Schrödinger discovered the
wave equation and contributed to the development of quantum mechanics. In 1927 Heisenberg proposed
the uncertainty principle according to which it is impossible to measure the precise values of momentum
and position of a subatomic particle. This way the modern quantum theory was developed in the early
20th century. As we have already seen, quantum physics mainly deals with waves and the subatomic
particles of matter. For this reason quantum theory is also referred to as quantum wave mechanics.
16.1 Heisenberg Uncertainty PrinciPLe LO1

Heisenberg uncertainty principle is perhaps the best known result of the wave particle duality, i.e., the concept
of waves or wave packet associated with a moving particle. According to Heisenberg uncertainty principle it
is impossible to determine simultaneously the exact position and momentum (or velocity) of a small moving
particle like electron.
As discussed earlier, the quantity |y(x, t)|2Dx represents the probability that the particle is within the region
between x and x + Dx. It means there is an uncertainty in the location of the position of the particle and Dx
is a measure of the uncertainty. The uncertainty in the position would be less if Dx is smaller, i.e., if the wave
packet is very narrow. The narrow wave packet means the range of wavelength Dl between l and l + Dl is
smaller or the range of wave numbers Dk between k and k + Dk is larger. So Dx is inversely proportional to
Dk, i.e.,
1
Dx µ
Dk
h h h 2p Dp
We may approximate this as DxDk ª 1. Taking  = , we get p = = = k , Dk = . Therefore
2p l 2p l 
DxDp ª 
The above relation represents the lowest limit of accuracy. Therefore, we can write more generally.
DxDp ≥ 
Dp
The principle of uncertainty can also be represented in terms of energy E and time t. Since = DF ,
Dt
we can write
Dp = DF ◊ Dt
Putting this value of Dp in the expression DxDp ≥  we obtain
Dx ¥ (DF ¥ Dt) ≥ 
or [DF ¥ Dx] Dt ≥ 
DE Dt ≥ 
The principle of uncertainty can also be expressed in terms of angular momentum and angle. Suppose we have
a particle at a particular angular position q and its angular momentum is Lq. Then the limits in the uncertainties
Dq and DLq are given by the relation DqDLq ≥ .
16.1.1 Mathematical Proof
Heisenberg’s uncertainty principle can be proved on the basis of deBroglie’s wave concept that a material
particle in motion is equivalent to a group of waves or wave packet, the group velocity G being equal to the
particle velocity v. Consider a simple case of wave packet which is formed by the superposition of two simple
harmonic plane waves of equal amplitudes a and having nearly equal frequencies w1 and w2. The two waves
can be represented by the equations.
y1 = a sin(w1t – k1x)
y2 = a sin(w2t – k2x)
w1 w
where k1 and k2 are their propagation constants and and 2 are their respective phase velocities. The
k1 k2
resultant wave due to superposition of these wave is given by
y = y1 + y2
È Dw Dk ˘ (i)
y = 2a sin (w t - kx) cos Í t- x˙
Î 2 2 ˚
where w = (w1 + w2)/2, k = (k1 + k2)/2, Dw = w1 – w2 and Dk = k1 – k2.
The resultant wave is shown in Fig. 16.1. The envelope (loop) of this wave travels with the group velocity G,
given by
Dw w1 - w 2
G= =
Dk k1 - k2
2a
Figure 16.1
Since the group velocity of deBroglie wave group associated with the moving particle is equal to the particle
velocity, the loop so formed is equivalent to the position of the particle. Then the particle may be anywhere
within the loop. Now the condition of the formation of node from Eq. (i) is given by
È Dw Dk ˘
cos Í t- x =0
Î 2 2 ˙˚
Dw Dk p 3p (2n + 1)p
or t- x= , , º, (ii)
2 2 2 2 2
where n = 0, 1, 2, …..
If x1 and x2 be the values of positions of two consecutive nodes, then from above equation by putting n and
(n + 1), we get
Dw Dk (2n + 1)p
t- x1 = (iii)
2 2 2
and Dw Dk (2n + 3)p (iv)

t- x2 =
2 2 2
From Eqs. (iii) and (iv), we have
Dk
( x1 - x2 ) = p
2
Dk
or Dx = p (v)
2
2p
or Dx =
Dk
2p 2p 2p p
but k= = =
l h/ p h
2p
Dk = Dp
h
where Dp is the error (uncertainty) in the measurement of momentum p. Therefore, from Eq. (v)
2p h h
Dx = =
2pDp Dp
or Dp Dx = h
However, more accurate measurements show that the product of uncertainties in momentum (Dp) and the
position (Dx) cannot be less than h/2p. Therefore
or DpDx ≥ 
This is the Heisenberg’s uncertainty principle.
16.1.2 Applications
Some important applications of uncertainty principle are discussed below.
16.1.2.1 Non-Existence of Electron in the Nucleus
The radius of the nucleus of an atom is of the order of 10–14 m. If an electron is confined within the nucleus,
the uncertainty in its position must not be greater than 10–14 m. According to uncertainty principle for the
lowest limit of accuracy
h
Dx Dp = (i)
2p
where Dx is uncertainty in the position and Dp is the uncertainty in the momentum.
From Eq. (i),
h 6.625 ¥ 10 -34
Dp = = (as Dx = diameter of nucleus)
2pDx 2 ¥ 3.14 ¥ 2 ¥ 10 -14
Dp = 5.275 ¥ 10 -21 kg m/sec
This is the uncertainty in momentum of the electron. It means the momentum of the electron would not be
less than Dp, rather it could be comparable to Dp. Thus
p = 5.275 ◊ 10–21 kg m/sec
The kinetic energy of the electron can be obtained in terms of momentum as
1 2 p2
T= mv =
2 2m
(5.275 ¥ 10-21 ) 2
= J
2 ¥ 9.1 ¥ 10-31
(5.275 ¥ 10-21 ) 2
= eV
2 ¥ 9.1 ¥ 10-31 ¥ 1.6 ¥ 10-19
= 95.55 ¥ 106 eV
ª 96 MeV
From the above result, it is clear that the electrons inside the nucleus may exist only when it possesses the
energy of the order of 96 MeV. However, the maximum possible kinetic energy of an electron emitted by
radioactive nuclei has been found about 4 MeV. Hence, it is concluded that the electron cannot reside inside
the nucleus.
16.1.2.2 Radius of Bohr’s First Orbit
If Dx and Dp be the uncertainties in determining the position and momentum of the electron in the first orbit,
then from the uncertainty principle
Dx Dp ª 

or Dp ª (i)
Dx
The uncertainty in kinetic energy (K.E.) of electron may be written as
( Dp ) 2 È p2 ˘ (ii)
DT = Í K.E. = T = ˙
2m Î 2m ˚
From Eqs. (i) and (ii), we have
2
1 È ˘
DT =
2m ÍÎ Dx ˙˚
and the uncertainty in the potential energy of the same electron is given by
1 ( Ze)( -e) È 1 ( Ze)( -e) ˘
DV = Í V = 4pe ˙
4pe 0 Dx Î 0 x ˚
The uncertainty in the total energy of electron together with Ze as the nucleus charge
DE = DT + DV
2 1 Ze 2
= -
2m( Dx) 2 4pe 0 Dx
The condition for this uncertainty in the energy to be minimum is
d ( DE )
=0
d ( Dx)
2 1 Ze 2
or - + =0
m( Dx)3 4pe 0 ( Dx) 2
 2 (4pe 0 )
Dx =
mZe2
Hence, the radius of first orbit

 2 (4pe 0 )
r = Dx =
mZe 2
e 0 h2
or r=
p mZe 2
This is the radius of first Bohr’s orbit.
16.1.2.3 Energy of a Particle in a Box or Infinite Potential Well
Let us consider a particle having mass m in infinite potential well of width L. The maximum uncertainty in
the position of the particle may be
(Dx)max = L
From the uncertainty principle
Dx Dp = 
 
or Dp = =
Dx L
Kinetic energy
p2 2
T= =
2m 2mL2
2
T=
2mL2
This is the minimum kinetic energy of the particle in an infinite potential well of width L.
16.1.2.4 Ground State Energy of Linear Harmonic Oscillator
The total energy E of a linear harmonic oscillator is the sum of its kinetic energy (K.E.) and potential
energy (P.E.).
E = K.E + P.E.
p2 1 2
E= + kx (i)
2m 2
Let a particle of mass m executes a simple harmonic motion along x-axis. The maximum uncertainty in the
determination of its position can be taken as Dx. From the uncertainty relation, the uncertainty in momentum
is then given by
 
Dp = [Taking Dp Dx = for more accuracy] (ii)
2 Dx 2
For maximum uncertainties Dp  p and Dx  x
Hence, the total energy E of the oscillator becomes
( Dp ) 2 1
E= + k ( Dx) 2
2m 2
2
Ê ˆ 1 1
or E=Á ˜ ◊ + k ( Dx) 2
Ë 2 ¯ 2m( Dx) 2 2
2 1
or E= 2
+ k ( Dx) 2 (iii)
8m( Dx) 2
For a minimum value of energy,

∂E
=0
∂( Dx)
Then we get
2
- + k ( Dx) = 0
4m( Dx)3
2
or ( Dx) 4 =
4mk
1/4
Ê 2 ˆ
( Dx) = Á
or Ë 4mk ˜¯ (iv)
Substituting value of Dx in Eq. (iii) from Eq. (iv), we get

1/2 1/2
 2 Ê 4mk ˆ 1 Ê 2 ˆ
Emin = Á ˜ + kÁ ˜
8m Ë  2 ¯ 2 Ë 4mk ¯
1/2
Ê k ˆ
Emin = Á ˜
2 Ë m¯
k
But = w = angular frequency. Therefore, the minimum energy of harmonic oscillator is expressed by
m
the following relation
1
Emin = w
2
Here it would be worth mentioning that the energy comes out to be Emin = w if we use DpDx =  (less
accuracy).
16.2 Wave FUnctiOn and its PHysicaL signiFicance

Waves in general are associated with quantities that vary periodically. In case of matter waves, the quantity
that varies periodically is called wave function. The wave function, represented by y, associated with the
matter waves has no direct physical significance. It is not an observable quantity. However, the value of the
wave function is related to the probability of finding the particle at a given place at a given time. The square
of the absolute magnitude of the wave function of a body evaluated at a particular time at a particular place
is proportional to the probability of finding the particle at that place at that instant.
The wave functions are usually complex. The probability in such a case is taken as y*y, i.e., the product
of the wave function with its complex conjugate, y* being the complex conjugate. Since the probability of
finding a particle somewhere is finite, we have the total probability over all space equal to unity. That is
•
Ú y *y dV = 1 (i)
-•
where dV = dxdydz.
Equation (i) is called the normalisation condition and a wave function that obeys this equation is said to be
normalised. Further, y must be a single valued since the probability can have only one value at a particular
place and time. Besides being normalisable, a further condition that y must obey is that it and its partial
∂y ∂y ∂y
derivatives , and be continuous everywhere.
∂x ∂y ∂z
The important characteristics of the wave function are as follows.
(i) y must be finite, continuous and single valued everywhere.
∂y ∂y ∂y
(ii) , and must be finite, continuous and single valued.
∂x ∂y ∂z
(iii) y must be normalisable.
16.3 TiMe indePendenT Schrödinger equATion LO2

Consider a system of stationary waves associated with a moving particle. The waves are said to be stationary
w.r.t the particle. If the position coordinates of the particle are (x, y, z) and y be the periodic displacement
for the matter waves at any instant of time t, then we can represent the motion of the wave by a differential
equation as follows.
∂ 2y ∂ 2y ∂ 2y 1 ∂ 2y (i)
2
+ 2
+ 2
=
∂x ∂y ∂z u 2 ∂t 2
where u is the velocity of wave associated with the particle. The solution of Eq. (i) gives y as a periodic
displacement in terms of time, i.e.,
y(x, y, z, t) = y (x, y, z)e–iwt (ii)
0
where y0 is the amplitude of the particle wave at the point (x, y, z) which is independent of time (t). It is a
function of (x, y, z). i.e., the position r and not of time t, Here.
r = xiˆ + yjˆ + zkˆ (iii)
Eq. (ii) may be expressed as
(iv)
y(r, t) = y0(r)e–iwt
Differentiating Eq. (iv) twice with respect to t, we get
∂ 2y
= - w 2y 0 (r )e - iw t
∂t 2
∂ 2y (v)
or 2
= - w 2y
∂t
∂ 2y
Substituting the value of from this equation in Eq. (i), we get
∂t 2
∂ 2y ∂ 2y ∂ 2y w2
+ + + y =0 (vi)
∂x 2 ∂y 2 ∂z 2 u2
where w = 2pn = 2p(u/l) [as u = ln]

so that
w 2p
= (vii)
u l
∂ 2y ∂ 2y ∂ 2y
Also 2
+ 2
+ 2
= —2y (viii)
∂x ∂y ∂z
where =2 is known as Laplacian operator. Using Eqs. (vi), (vii) and (viii), we have
4p 2 (ix)
—2y + 2
y =0
l
Also from the deBroglie wave concept
h
l=
mv
Using this relation in Eq. (ix) gives
4p 2 m 2 v 2 (x)
—2y + 2
y =0
h
Here it can be noted that the velocity of particle v has been introduced in the wave equation.
If E and V are respectively the total energy and potential energy of the particle then its kinetic energy is
given by
1 2
mv = E - V
2
m 2 v 2 = 2 m( E - V ) (xi)
The use of Eq. (xi) in Eq. (x) gives rise to
8p 2 m
—2y + ( E - V )y = 0
h2
or 2m (xii)
—2y + ( E - V )y = 0
2
This is the time independent Schrödinger equation, where the quantity y is known as wave function.
For a freely moving or free particle V = 0. Therefore, Eq. (xii) becomes
2mE
—2y + y =0 (xiii)
2
This is called time independent Schrödinger equation for a free particle.
16.4 TiMe dePendenT Schrödinger equATion LO3

In order to obtain a time dependent Schrödinger equation, we eliminate the total energy E from time
independent Schrödinger equation. For this we differentiate Eq. (iv) w.r.t. t and obtain
∂y
= - iwy 0 (r )e - iw t
∂t
= i (2pn )y 0 (r )e - iw t
E iE i
= - 2pn iy = - 2p i y =- ¥ y
h h i
∂y Ey
fi =
∂t i
∂y
or Ey = i (xiv)
∂t
Substituting the value of Ey from Eq. (xiv) in Eq. (xii), we have
2 m È ∂y ˘
—2y + 2 Í
i - Vy ˙ = 0
 Î ∂ t ˚
2m È ∂y ˘
or —2y = - 2 Í
i - Vy ˙
 Î ∂t ˚
(xv)
Ê 2 2 ˆ ∂y
or ÁË - — + V ˜ y = i
2m ¯ ∂t
Ê 2 2 ˆ
This equation is known as Schrödinger’s time dependent wave equation. The operator Á - — + V ˜ is
Ë 2m ¯
called Hamilitonian operator and is represented by H. If we see the RHS of Eq. (xv) and keep in mind
∂
Eq. (xiv), we notice that the operator i operating on y gives E. Hence, Schrödinger equation can be
written in operator form, as below ∂ t
Hy = Ey
16.5 OPeratOrs LO3

In a physical system, there is a quantum mechanical operator that is associated with each measurable parameter.
In quantum mechanics, we deal with waves (wave function) rather than discrete particles whose motion and
dynamics can be described with the deterministic equations of Newtonian physics. Generally an operator
is anything that is capable of doing something to a function. There is an operator corresponding to every
observable quantity. However, the choice of operator is arbitrary in quantum mechanics. When an operator
operates on a wave function it must give observable quantity times the wave function. It is a must condition
for an operator.
If we consider an operator represented by A corresponding to the observable quantity a, then
Ay = ay
Wave function that satisfies the above equation is called eigen function and corresponding observable
quantity is called eigen value and the equation is called eigen value equation. Some of those operators are
tabled below.
Classical Quantity Quantum Mechanical Operator

Position x, y, z, x, y, z

Momentum p –i—
Momentum
components px, py, pz ∂ ∂ ∂
-i  , -i , -i
∂x ∂y ∂z
Energy E ∂
i
∂t
Hamiltonian 2 ∂2
- + V (r )
(Time independent) 2 m ∂x 2
2 ∂2
Kinetic energy -
2m ∂r 2
16.6 APPlicATionS of Schrödinger equATion LO4

 
In classical mechanics, based on Newton’s second law of motion ( F = ma ) we make a mathematical
prediction of the path, a given system will take following set of known initial conditions. The analogue of
Newton’s law is Schrödinger equation in quantum mechanics for a quantum system such as atoms, molecules
and subatomic particles. The subatomic particles may be free, bound or localized. Schrödinger equation
describes the time evolution of the system’s wave function.
Schrödinger’s equation is extremely useful for investigating various quantum mechanical problems. With
the help of this equation and boundary conditions, the expression for the wave function is obtained. Then the
probability of finding the particle is calculated by using the wave function. In the following subsections, we
discuss different quantum mechanical problems, viz. particle in a box, one-dimensional harmonic oscillator,
step potential and step barrier.
16.6.1 Particle in a Box (infinite Potential Well)
The simplest quantum mechanical problem is that of a particle trapped in a box with infinitely hard walls. Infinitely
hard walls means the particle does not loose energy when it collides with such walls, i.e., its total energy remains
constant. A physical example of this problem could be a molecule which is strictly confined in a box.
Let us consider a particle restricted to move along the x-axis between x = 0 and x = L,
by ideally reflecting, infinitely high walls of the infinite potential well, as shown in V
Fig. 16.2. Suppose that the potential energy V of the particle is zero inside the box, but
rises to infinity outside, that is,
0 X
L
V = 0 for 0 £ x £ L
V = • for x < 0 and x > L Figure 16.2
In such a case, the particle is said to be moving in an infinitely deep potential well. In order to evaluate the wave
function y in the potential well, Schoredinger equation for the particle within the well (V = 0) is written as
∂ 2y 8p 2 mE
+ y=0 (i)
∂x 2 h2
8p 2 mE
we put = k 2 in the above equation for getting
h2
∂ 2y
+ k 2y = 0 (ii)
∂x 2
The general solution of this differential equation is
y(x) = A sin kx + B cos kx (iii)
where A and B are constants.
Applying the boundary condition y(x) = 0 at x = 0, which means the probability of finding particle at the wall
x = 0 is zero, we obtain
A sin (0) + B cos (0) = 0 fi B=0
Again, we have y(x) = 0 at x = L, then
A sin kL + B cos kL = 0 fi A sin kL = 0
The above equation is satisfied when
kL = np
np
or k= where n = 1, 2, 3, …
L
or n 2p 2
k2 = (iv)
L2
8p 2 mE n 2p 2
or = (v)
h2 L2
or in general we can write Eq. (v) as
n2 h2
En = where n = 1, 2, 3, …
8mL2
Thus, it can be concluded that in an infinite potential well the particle cannot have
an arbitrary energy, but can take only certain discrete energy values corresponding
to n = 1, 2, 3, …. These are called the eigen values of the particle in the well and
constitutes the energy levels of the system. The integer n corresponding to the energy
level En is called its quantum number, as shown in Fig. 16.3.
We can also calculate the momentum p of the particle or the eigen values of the
momentum, as follows,
2p 2p p
Since k = = =
l h/ p  Figure 16.3
np 
p = k =
L
The wave function (or eigen function) is given by Eq. (iii) along with the use of expression for k.
np x
y n ( x) = A sin
L
To find the value of A, we use the normalisation condition.

•
Ú | y n ( x) |
2
dx = 1
-•
As mentioned earlier, the above expression simply says that the probability of finding the particle is 1. In the
present case, the particle is within the box i.e., between 0 < x < L. So the normalisation condition becomes
L
np x
A2 Ú sin 2 dx = 1
0
L
Ê Lˆ 2
A2 Á ˜ = 1 or A=
Ë 2¯ L
The normalised eigen wave function of the particle is, therefore, given by
2 np x
y n ( x) = sin
L L
The first three eigen functions y1, y2, y3 together with the probability densities |y1|2, |y2|2,|y3|2, are shown
in Figs. 16.4(a) and (b), respectively.
Figure 16.4
Classical mechanics predicts the same probability for the particle being anywhere in the well. Wave
mechanics, on the other hand, predicts that the probability is different at different points and there are points
(nodes) where the particle is never found. Further, at a particular point, the probability of finding the particle
is different for different energy states. For example, a particle in the lowest energy state (n = 1) is more likely
to be in the middle of the box, while in the next energy state (n = 2) it is never there since |y2|2 is zero there.
It is |yn|2 which provides the probability of finding the particle within the potential well.
16.6.2 finite Potential Step

A physical example of this quantum mechanical problem can be thought as the neutron which is trying to
escape nucleus. The potential function of a potential step may be represented as
V(x) = 0 for x < 0 region I
V(x) = V0 for x > 0 region II (i)
We consider that a particle of energy E is incident from left on the potential step of height V0 as shown in
Fig. 16.5. Further, we assume that the energy of the incident particle is greater than the step barrier height
i.e., E > V0. Since E > V0, according to classical theory there should be no reflection at the boundary of
the step potential barrier. However, quantum mechanically this is not true. It means that there will be some
reflection from the boundary of the potential step.
V(x)
Region II
Region I ⇐ Potential and Energy
E
VO
X
∗
y y
⇐ Probability Density
X
O
Figure 16.5
The wavelength of the particle suddenly changes from region I to region II and is given as follows
h h
l1 = = (ii)
p1 2mE
h h
and l2 = = (iii)
p2 2m( E - V0 )
Hence, a small part of the wave associated with the particle is reflected due to this change in wavelength
and the rest part is transmitted. This can be proved with the solution of Schrödinger wave equations for two
regions. The Schrödinger equation for region I is written as
d 2y 1 ( x) 2mE
+ y 1 ( x) = 0 (iv)
dx 2 2
Schrödinger equation for region II is written as
d 2y 2 ( x) 2m( E - V0 )
+ y 2 ( x) = 0 (v)
2
dx 2
The solutions of Eqs. (iv) and (v) are written as
y 1 ( x) = A1eik1x + A2 e - ik1x (vi)
y 2 ( x) = A3eik2 x + A4 e - ik2 x (vii)
where y1(x) and y2(x) are the wave functions of region I and II and A1, A2, A3 and A4 are constants. k1 and k2
are defined as follows
2mE 2m( E - V0 )
k1 = and k2 =
 
The first term in Eq. (vi) represents the wave travelling in the positive x direction in the first region and
second term represents the reflected part of the incident wave travelling in the negative x direction in region I.
In Eq. (vii), first term represents the transmitted part of the incident particle wave travelling in the direction of
positive x axis in region II. The second term of Eq. (vii) has no meaning, because the reflection of the particle
cannot take place in region II. So, considering this, Eq. (vii) can be written as
y 2 ( x) = A3eik2 x (viii)
The boundary condition at x = 0 is defined as
y1(0) = y2(0), (ix)
which means the wave function is continuous at the boundary. Also the derivative of y should be continuous
at the boundary, i.e.,
dy 1 ( x ) dy 2 ( x ) (x)
=
dx x=0 dx x=0
Using these boundary conditions, we get

A1 + A2 = A3 (xi)
and ik1 ( A1eik1x - A2 e - ik1x ) x = 0 = ik2 A3eik2 x x=0
or ik1 ( A1 - A2 ) = ik2 A3 (xii)
Solving Eq. (xi) and (xii), we get
A2 k1 - k2
=
A1 k1 + k2 (xiii)
A3 2k1
and = (xiv)
A1 k1 + k2
Since the reflection and transmission of the particle takes place, the problem can be investigated based on the
reflection and transmission coefficients. Further, the coefficient A1 is related to the wave function y1, i.e., of
the incident particle and A3 is related to the wave function y3, i.e., of the transmitted particle. It means the
reflection and transmission coefficients can be defined as follows.
Reflection coefficient = Reflected intensity/Incident intensity
= (Reflected amplitude)2/(Incident amplitude)2 (xv)
Transmission coefficient = Transmitted intensity/Incident intensity
= (Transmitted amplitude)2/(Incident amplitude)2 (xvi)
The reflection coefficient is
2
| A2 |2 Ê k1 - k2 ˆ (xvii)
R= = π0
| A1|2 ÁË k1 + k2 ˜¯
and transmission coefficient is

2
| A3 |2 Ê 2k1 ˆ
T= = π1 (xviii)
| A1|2 ÁË k1 + k2 ˜¯
From the above results we see that the reflection coefficient (R) is not zero and the transmission probability
(T) is not unity in the quantum mechanical treatment of the particle behaviour in the finite potential step
problem. However, classically the reflection coefficient should be zero and transmission coefficient should
be equal to unity.
16.6.3 finite Potential Barrier

A physical example of this quantum mechanical problem can be thought as the a particle which is trying to
escape Coulomb barrier. For this case, the potential function is defined as
V(x) = 0 for x < 0 Region I
V(x) = V0 for 0 < x < a Region II
V(x) = 0 for x > a Region III
The potential barrier is considered between x = 0 and x = a, as shown in Fig. 16.6. Here, we suppose that a
particle incident on the barrier has energy E which is less than the barrier height V0, i.e., E < V0. Classically,
when E < V0, the particle can never penetrate the potential barrier and appear in region III. It means the
particle is always reflected from the barrier. Therefore, the transmission coefficient is zero. However, quantum
mechanically this is not true and there is some probability for a particle penetrate the barrier. It means a
fraction of the particles incident from the left will cross the barrier and appear in region III.
V(x)
VO
⇐ Potential and Energy
E Region I Region II Region III
X
0 a
⇐ Probability Density
y∗y
X
0 a
Figure 16.6
Schrödinger wave equations for the regions I and III are as follows
d 2y 1 ( x) 2mE
+ y 1 ( x) = 0 (ii)
2
dx 2
d 2y 3 ( x) 2mE (iii)
and 2
+ 2
y 3 ( x) = 0
dx 
where y1 (x) and y3 (x) are the wave function of region I and III. The solutions of these equations are
y 1 ( x) = A1eik1x + A2 e - ik1x (iv)
and y 3 ( x) = A3eik1x + A4 e - ik1x (v)

2mE p 2p
where k1 = = = and A1, A2, A3 and A4 are constants. The solution for y1 is a combination of
  l
reflected and transmitted wave in region I. But in the region III, the reflected part of the wave is zero (A4 =
0) and the transmitted wave is traveling in the positive x direction. So the solution in region III becomes
y 3 ( x) = A3eik1x (vi)
Now, the Schrödinger equation for region II is written as
d 2y 2 ( x) 2m( E - V0 )
2
+ y 2 ( x) = 0 (vii)
dx 2
But as we know that E < V0 then it will be convenient to write this equation in the form
d 2y 2 ( x) 2m(V0 - E )
2
- y 2 ( x) = 0 (viii)
dx 2
where y2 is the wave function of region II. The solution of above equation is
y 2 ( x) = A5 e - ik2 x + A6 eik2 x (ix)
2m(V0 - E )
where k2 =

In order to calculate the transmission probability T, we must apply the boundary conditions to wave function
y1, y2 and y3. These boundary conditions at the left hand (at x = 0) or at the right hand wall (at x = a) of the
barrier are defined as
Boundary conditions at x = 0 are
y1(0) = y2(0) (x)
∂y 1 (0) ∂y 2 (0)
and = (xi)
∂x ∂x
At x = a are
y2(a) = y3(a) (xii)
∂y 2 ( a ) ∂y 3 ( a ) (xiii)
and =
∂x ∂x
The above boundary conditions along with the use of wave functions y1, y2 and y3 yield
A1 + A2 = A5 + A6 (xiv)
ik1A1 – ik1A2 = –k2A5 + k2A6 (xv)
- k2 a k2 a ik1a (xvi)
A5 e + A6 e = A3e
- k2 A5 e - k2 a + k2 A6 e k2 a = ik1 A3eik1a (xvii)
Solving Eqs. (xiv) and (xv), we get

(ik1 - k2 ) (ik + k )
A1 = A5 + 1 2 A6 (xviii)
2ik1 2ik1
and solving Eqs. (xvi) and (xvii), we get

(k2 + ik1 )eik1a
A6 = A3 (xix)
2 k 2 e k2 a
(k2 - ik1 )eik1a (xx)

and A5 = - k2 a
A3
2k 2 e
Substituting these values in Eq. (xviii), we get
A1 È 1 i Ê k2 k1 ˆ ˘ (ik1 + k2 ) a È 1 i Ê k2 k1 ˆ ˘ (ik1 - k2 ) a
= + - e + Í - Á - ˜˙e (xxi)
A3 ÍÎ 2 4 ÁË k1 k2 ˜¯ ˙˚ Î 2 4 Ë k1 k2 ¯ ˚
As we assumed, the potential barrier to be much higher than the energy of the incident particles. In this
condition k2/k1 >> k1/k2 and hence
k2 k1 k2
- ª (xxii)
k1 k2 k1
Further we assume that the potential barrier wide enough so that for y2 gets severely weakened between x =
0 and x = a. This means k2a >> 1, i.e.,
e k2 a >> e - k2 a (xxiii)
So Eq. (xxi) is approximated by
A1 Ê 1 ik2 ˆ (ik1 + k2 ) a
= + e (xxiv)
A3 ÁË 2 4k1 ˜¯
The complex conjugate of this is written as

*
Ê A1 ˆ Ê 1 ik2 ˆ ( - ik1 + k2 ) a
ÁË A ˜¯ = ÁË 2 - 4k ˜¯ e (xxv)
3 1
On multiplying Eqs. (xxiv) and (xxv), we get
A1 A1* Ê1 k22 ˆ 2 k2 a
So, = Á + ˜e (xxvi)
A3 A3* Ë 4 16k12 ¯
Since the coefficient A1 is related to the wave function y1, i.e., of the incident particle and A3 is related to the
wavelength of y3, of the transmitted particle, the transmission probability is equivalent to
-1
A A* Ê A A* ˆ Ê 16 ˆ -2 k2a
T = 3 3* = Á 1 1* ˜ =Á e (xxvii)
A1 A1 Ë A3 A3 ¯ Ë 4 + (k2 /k1) 2 ˜¯
From the definitions of k1 and k2 we see that

2
Ê k2 ˆ 2m(V0 - E )/  2 V0
ÁË k ˜¯ = = -1 (xxviii)
1 2mE /  2 E
With this it can be seen that the quantity in the bracket varies slowly with E and V0 than the variation of
exponential term. So the approximated transmission probability is
T = e - k2 a (xxix)
16.6.4 one-dimensional harmonic oscillator
A physical example of this quantum mechanical problem can be thought as an atom of vibrating diatomic
molecule. In general, a particle undergoing simple harmonic motion in one dimension is called one
dimensional harmonic oscillator. The potential and total energy of such a system is shown in Fig. 16.7 where
the probability density is also shown. In such a motion, the restoring force F is proportional to the particle’s
displacement x from the equilibrium position, i.e.,
V(x)
F = –kx (i)
where k is force constant. The potential energy V can be written as E
1 2
V= kx
2
1 X
Then, the Schrödinger’s equation for the oscillator with V = kx 2 is
2
d 2y 2m È 1 2˘
+ 2 Í E - kx ˙ y = 0
dx 2  Î 2 ˚
y�y
1/2
h 8p mE 2 Ê 4p mk ˆ 2
X
Putting = , = a and Á =b in the above
2p h 2 Ë h 2 ˜¯ Figure 16.7
equation, we obtain
d 2y
+ (a - b 2 x 2 )y = 0 (ii)
dx 2
Now we introduce a dimension less independent variable as x = b x . Thus Eq. (ii) becomes.
d 2y È x2 ˘
b 2
+ Ía - b 2 ˙ y = 0
dx Î b˚
d 2y È a ˘
+ - x 2 ˙y = 0 (iii)
d x 2 ÍÎ b ˚
The solution of this equation is
2
y = CUe -x /2
(iv)
where U is a function of x. Then Eq. (iii) takes the form
d 2U dU È a ˘
2
- 2x + Í - 1˙ U = 0
dx dx Î b ˚
a
If we replace - 1 by 2n, this equation becomes Hermite differential equation. Then function U(x) may be
b
replaced with Hermite polynomial H. So, we get
d 2H dH
2
- 2x + 2nH = 0
dx dx
Thus, the solution of Eq. (iii) is obtained by replacing U by Hermite polynomial H in Eq. (iv). Hence, we get
2
y = CHe -x /2
2
In general, y n (x ) = CH n (x )e -x /2
, where n = 0, 1, 2,…
eigen Values of energy

a
Since - 1 = 2n
b
a 8p 2 mE 4p 2 mk
fi = 2n + 1 fi a = (2n + 1)b fi = (2n + 1)
b h2 h2
This restriction gives a corresponding restriction on E, i.e.,
Ê 1ˆ h k
E = Án + ˜
Ë 2 ¯ 2p m
1 k
But = n is the frequency of oscillations. Hence, the energy can be written in terms of n as
2p m
Ê 1ˆ
E = Á n + ˜ hn
Ë 2¯
n En
Thus, in general, the oscillator has finite, unambiguous and continuous
solutions at values of E given by
Ê 1ˆ (v)
En = Á n + ˜ hn 4 E4 = 9 hn
Ë 2¯ 2
Following conclusions can be drawn from equation (v) 3 E3 = 7 hn

2
The particle executing simple harmonic motion can have only
2 E2 = 5 hn
2
(i) discrete energy levels that are equidistant and are separated by
hn, as shown in Fig. 16.8 1 E1 = 3 hn
2
(ii) The energy levels are non-degenerate.
0 E0 = 1 hn
2
1
(iii) For n = 0, E0 = hn . It means the minimum energy is not
zero. 2 Figure 16.8
16.7 quAnTuM STATiSTicS LO5

Classical statistics, i.e., Maxwell-Boltzmann statistics, successfully explained the energy and velocity
distribution of molecules of an ideal gas but it failed to explain the energy distribution of electrons in metals, for
example, electron gas and the energy distribution of photons in a photon gas. These phenomena can be explained
on the basis of quantum statistics, where the particles of the system are considered to be indistinguishable
contrary to the consideration of particles as distinguishable in classical statistics. If ni particles are distributed
in the gi cell of the ith compartment in the phase space, then the number of particles per cell is defined as
ni(E)/gi(E). The factor ni(E)/gi(E) is called occupation index. If ni(E)/gi(E) ≥ 0 or 1 the particles are considered
as indistinguishable which is the basic feature of the quantum statistics. If the indistinguishable particles have
integral spin, we use the Bose-Einstein distribution function and if the particles have half-integral spins, then
Fermi-Dirac distribution function is appropriate. The brief description of these statistics is given below.
16.7.1 Bose-einstein Statistics
It is applicable to those systems which contain identical, indistinguishable particles of zero or integral spins.
Such particles are called bosons. Examples of bosons are photons, phonons etc. Pauli exclusion principle
does not apply to the bosons. Bose-Einstein distribution law is given by
gi ( E )
ni ( E ) = a + bE (i) ni (E) / gi (E)
e -1
8
where a = - EF /kT and b = 1/kT. This law is also applicable in the
case of photon gas for which b = 0 and E = hn. For the photon gas,
the above equation reads 6
gi ( E ) (ii)
ni ( E ) = T2 > T1
eb E - 1 4
The plot of ni(E)/gi(E) versus E is shown in Fig. 16.9 for two

different temperatures with T2 > T1. If E >> kT, the exponential term
2 T2
in the above equation is very large and –1 may be dropped. It means
T1
ni(E) = gi(E)e–E/kT
E
The above relation represents Maxwell-Boltzmann statistics. So 0 1 2 3
the Bose-Einstein statistics reduces to the Maxwell-Boltzmann Figure 16.9
statistics under the condition E >> kT.
At low energy, i.e., when E << kT, ebE can be neglected as –1 predominates. This makes ni(E)/gi(E) much
larger for Bose-Einstein statistics than for Maxwell-Boltzmann statistics at low energies.
16.7.2 fermi-dirac Statistics
This statistics is applicable to systems, which consist of identical, independent and indistinguishable particles of
having half-integral spins. The particles, which obey Fermi-Dirac statistics, are called fermions. The examples
of fermions are electrons, protons, neutrons, etc. The fermion must obey Pauli exclusion principle. In Fermi-
Dirac statistics, interchange of two particles of the system leaves the resultant system in an antisymmetric state.
That is, the wave function of the system gets changed only with minus sign. As it obeys the Pauli exclusion
principle, in Fermi-Dirac statistics, there can be only one particle in each state. Hence, the total number of
particles must be less than or equal to the total number of states available. Under these considerations, fermions
lead to the following distribution law, named Fermi-Dirac distribution law, given by
gi ( E )
ni ( E ) = a +bE (i)
e +1
where a = –EF/kT and b = 1/kT.

gi ( E )
So, ni ( E ) = ( E - EF )/ kT (ii)
e +1
ni ( E ) 1
= ( E - E )/ kT (iii)
gi ( E ) e F +1
In the above equations, ni is the number of particles in an energy state E, gi is the statistical weight factor and
EF is the Fermi energy. Fermi energy is independent of temperature. The plots of ni(E)/gi(E) versus E for
different temperatures is shown in Fig. 16.10.
ni (E)
gi (E)
T= 0 T2 > T1
1
T1
T2
0 E
EF
Figure 16.10
Now the different cases will be discussed.

(a) At T = 0 K and when E < EF, Eq. (iii)
ni ( E )
=1 (iv) ni (E)
gi ( E ) gi (E)
It means that ni(E) = gi(E), i.e., all the energy states will have one
electron each. T= 0
(b) At T = 0 K and E > EF, 1
ni ( E )
=0 (v)
gi ( E )
It means that ni= 0, i.e., all such energy states which have
energies greater than Fermi energy are vacant. This clarifies that
all states with energies up to EF are filled while all states with 0 E
EF
energy greater than EF are vacant. The plot between ni(E)/gi(E)
Figure 16.11
and E is shown in Fig. 16.11 for these conditions at T = 0 K.
s UMMARY

✦ Quantum mechanics deals with the behaviour and characteristics of matter, in the subatomic level, and
energy.
✦ Basics of quantum theory were particularly developed by Planck, Einstein, Schrödinger and Heisenberg.
✦ The best known result of the wave particle duality, i.e., the concept of waves or wave packet associated
with a moving particle, is the Heisenberg uncertainty principle. According to this principle it is
impossible to determine simultaneously the exact position and momentum (or velocity) of a small
moving particle like electron. If Dx be the uncertainty in the position and Dp be the uncertainty in the
momentum, then DxDp ª  represents the lowest limit of accuracy. More generally, the Heisenberg
uncertainty principle in mathematical form is written as DxDp ª .
✦ Based on Heisenberg uncertainty principle (i) it was proved that the electron cannot reside inside the
nucleus, (ii) radius of Bohr’s first orbit was calculated, (iii) energy of a particle in a box was evaluated,
and (iv) ground state energy of a linear harmonic oscillator was calculated.
✦ Waves in general are associated with quantities that vary periodically. In case of matter waves, the
quantity that varies periodically is called wave function (y). The wave function y is generally a
complex quantity and as such it has no physical meaning. However, the product yy*, where y* is its
complex conjugate, represents the probability.
✦ Wave function y should be finite, continuous and single valued and so it is space derivative. Also, the
wave function should be normalisable.
2m
✦ Time independent Schrödinger equation is written as — 2y + ( E - V )y = 0, where y is the wave
2
function, m is the mass of the particle, E is the total energy of the particle and V is the potential.
2 2
✦ Time dependent Schrödinger equation is written as Hy = Ey, where H = — + V is Hamiltonian
2m
operator. E is the total energy of the particle.
n2 h2 ,
✦ Particle in a box acquires only discrete energy values, the energy being given by En =
8mL2
where m is the mass of the particle, L is the width of the box (infinite potential well) and
n (=1, 2, 3,…) is called the eigen value or the quantum number. It is interesting to note that for n = 1,
the particle is mostly found in the middle of the well, whereas the probability of finding the particle in
middle for n = 2 is zero. It happens quantum mechanically.
✦ In the quantum mechanical problem of finite potential well, when the particle energy E is less than the
barrier height V0, we find that there is a probability that the particle would cross the barrier. Here the
-k a 2m(V0 - E )
transmission probability was obtained as T = e 2 , where k2 = .

✦ In quantum statistics, the particles of the system are considered to be indistinguishable contrary to
the consideration of particles as distinguishable in classical statistics. If ni particles are distributed in
the gi cell of the ith compartment in the phase space, then the number of particles per cell is defined
as ni(E)/gi(E), which is called occupation index. If ni(E)/gi(E) ≥ 0 or 1, the particles are considered as
indistinguishable which is the basic feature of the quantum statistics.
✦ Quantum statistics has two branches, namely Bose-Einstein statistics and Fermi-Dirac statistics.
✦ Bose-Einstein statistics is applicable to those systems which contain identical, indistinguishable
particles of zero or integral spins. Such particles are called bosons. Examples of bosons are photons,
phonons etc. Pauli exclusion principle does not apply to the bosons.
✦ Fermi-Dirac statistics is applicable to systems, which consist of identical, independent and
indistinguishable particles having half-integral spins. The particles, which obey Fermi-Dirac statistics,
are called fermions. The examples of fermions are electrons, protons, neutrons, etc. The fermion must
obey Pauli exclusion principle. In Fermi-Dirac statistics, interchange of two particles of the system
leaves the resultant system in an antisymmetric state. That is, the wave function of the system gets
changed only with minus sign.
s olved e XAMPles
E xamplE 1 The position and momentum of a 1.0 keV electron are simultaneously measured. If the position
is located within 1 Å, what is the percentage of uncertainty in momentum?
Solution Given Dx = 1.0 ¥ 10–10 m and E = 1000 ¥ 1.6 ¥ 10–19 J =1.6 ¥ 10–16 J.
Heisenberg’s uncertainty principle says

DxDp = and p = 2mE
2
p = 2 ¥ 9.1 ¥ 10 -31 ¥ 1.6 ¥ 10 -16
= 1.71 ¥ 10 -23 kg m/sec
and
 h 6.62 ¥ 10 -34
Dp = = =
2 Dx 2 ¥ 2p ¥ Dx 2 ¥ 2 ¥ 3.14 ¥ 1.0 ¥ 10 -10
= 5.27 ¥ 10 -25 km g/sec
Percentage of uncertainty in momentum

Dp 5.27 ¥ 10 -25
= ¥ 100 = ¥ 100
p 1.71 ¥ 10 -23
= 3.1%
E xamplE 2 The uncertainty in the location of a particle is equal to its deBroglie wavelength. Calculate the
uncertainty in its velocity.
Solution Given Dx = h .
p
Now
 h
DxDp = =
2 4p
h 1 h p mv
or Dp = D (mv) = = =
4p Dx 4p h 4p
mv
mDv =
4p
v
or Dv =
4
E xamplE 3 The position and momentum of 0.5 keV electron are simultaneously determined. If its position
is located within 0.2 nm, what is the percentage uncertainty in its momentum?
Solution Given E = 0.5 ¥ 103 ¥ 1.6 ¥ 10–19 = 0.8 ¥ 10–16 J and Dx = 0.2 ¥ 10–9 m.
Now

DxDp = and momentum p = 2mE
2
so p = 2 ¥ 9.1 ¥ 10 -31 ¥ 0.8 ¥ 10 -16 = 12.06 ¥ 10 -24
or p = 1.21 ¥ 10 -23 kg m/sec
 1 h 1
or Dp = =
2 Dx 4p 0.2 ¥ 10 -9
6.62 ¥ 10 -34
Dp = = 2.635 ¥ 10 -25 kg m/sec
4 ¥ 3.14 ¥ 0.2 ¥ 10 -9
\ Percentage uncertainty in momentum
Dp 2.635¥ 10 -25
¥ 100 = ¥ 100
p 1.21 ¥ 10 -23
2.635 ¥ 10 -23
= = 2.18%
1.21 ¥ 10 -23
E xamplE 4 Wavelengths can be determined with accuracies of one part in 106. What is the uncertainty in the
position of a 1 Å X-ray photon when its wavelength is simultaneously measured?
Solution Given l = 10–10 m.
By uncertainty principle,
 h
DxDp = =
2 4p
h
and l= or pl = h (i)
p
By differentiating
pDl + lDp = 0
pDl hDl È h˘
or Dp = = 2 ÍÎ p = l ˙˚ (ii)
l l
By using Eqs. (i) and (ii), we get
hDl h
Dx =
l2 4p
l2
or DxDl = (iii)
4p
Wavelength can be measured with accuracy of one part in 106, it means the uncertainty in wavelength is
Dl 1
= = 10 -6 (iv)
l 106
By putting this value in Eq. (iii), then
Dl l l
Dx = or Dx ¥ 10 -6 =
l 4p 4p
106 ¥ l 106 ¥ 10 -10
or Dx = = = 7.96 mm
4p 4 ¥ 3.14
E xamplE 5 Calculate the uncertainty in measurement of momentum of an electron if the uncertainty in

locating it is 1 Å.
Solution Given Dx = 1.0 ¥ 10–10 m.
Formula used is
 h
DxDp = =
2 4p
h 1 6.62 ¥ 10 -34 1
Dp = = ¥ -10
4p Dx 4 ¥ 3.14 10
Dp = 5.27 ¥ 10 -25 kgm/sec
E xamplE 6 An electron has a momentum 5.4 × 10–26 kg m/sec with an accuracy of 0.05%. Find the minimum
uncertainty in the location of the electron.
Solution Given p = 5.4 × 10–26 kg m/sec.
The uncertainty in the measurement of momentum
5.4 ¥ 10 -26 ¥ 0.05
D=
100
= 2.7 ¥ 10 -29 kg m/sec
 h
DxDp = =
2 4p
h 1 6.62 ¥ 10 -34 1
\ Dx = = ¥
4p Dp 4 ¥ 3.14 2.7 ¥ 10 -29
= 1.952 ¥ 10 -6 m
= 1.952 mm
E xamplE 7 A hydrogen atom is 0.53 Å in radius. Use uncertainty principle to estimate the minimum energy
an electron can have in this atom.
Solution Given Dxmax = 0.53 Å.

Heisenberg’s uncertainty principle
 h
DxDp = =
2 4p
h
( Dx) max ( Dp ) min = (i)
4p
2
pmin ( Dp ) 2min
and (K.E.) min = = [ pmin = Dpmin ]
2m 2m
h 1 6.62 ¥ 10-34 1
( Dp ) min = =
4p Dx 4 ¥ 3.14 0.53 ¥ 10-10
= 0.9945 ¥ 10-24
= 9.945 ¥ 10-25 kg m/sec
( Dp ) 2min (9.945 ¥ 10-25 ) 2
and (K.E.) min = =
2m 2 ¥ 9.1 ¥ 10-31
= 5.434 ¥ 10-17 J
E xamplE 8 The speed of an electron is measured to be 5.0 × 103 m/sec to an accuracy of 0.003%. Find the
uncertainty in determining the position of this electron.
Solution Given v = 5.0 × 103 m/sec.
Formula used is
 h
DxDp = =
2 4p
0.003 0.003
Dv = v ¥ = 5.0 ¥ 103 ¥ = 0.15 m/sec
100 100
and Dp = mDv = 9.1 ¥ 10-31 ¥ 0.15 = 1.365 ¥ 10-31 kg m/sec
6.62 ¥ 10-34 1 6.62 ¥ 10-34 1
Dx = =
4 ¥ 3.14 Dp 4 ¥ 3.14 1.365 ¥ 10-31
= 3.861 ¥ 10-4 m
E xamplE 9 An electron has speed of 6.6 × 104 m/sec with an accuracy of 0.01%. Calculate the uncertainty in
position of an electron. Given mass of an electron as 9.1 × 10–31 kg and Planck’s constant h as 6.6 × 10–34 J sec.
0.01
Solution Given v = 6.6 ¥ 104 m/sec and Dv = 6.6 ¥ 104 ¥ m/sec
100
= 6.6 m/sec.
Formula used is
 h h 1
DxDp = = or Dx =
2 4p 4p Dp
Dp = mDv = 9.1 ¥ 10 -31 ¥ 6.6
h 1 6.6 ¥ 10 -34
or Dx = =
4p Dp 4 ¥ 3.14 ¥ 9.1 ¥ 10 -31 ¥ 6.6
Dx = 8.75 ¥ 10 -6 m
E xamplE 10 Calculate the smallest possible uncertainty in the position of an electron moving with a velocity
3 ¥ 107 m/sec.
Solution Given v = 3 ¥ 107 m/sec.
Formula used is
 h
DxDp = =
2 4p
m0v
Dpmin ª p = mv =
1 - v 2 /c 2
h 1 h È 1 - v 2 /c 2 ˘
\ Dx = = Í ˙
4p Dp 4p ÎÍ m0v ˚˙
È 7 2 ˘
Í 1 - Ê 3 ¥ 10 ˆ ˙
Á ˜
6.62 ¥ 10 -34 Í Ë 3 ¥ 108 ¯ ˙
= Í ˙
4 ¥ 3.14 ÍÎ 9.1 ¥ 10 -31 ¥ 3 ¥ 107 ˙˚
= 1.92 ¥ 10 -12 m
E xamplE 11 If an excited state of hydrogen atom has a life-time of 2.5 ¥ 10–14 sec, what is the minimum
error with which the energy of this state can be measured? Given h = 6.62 ¥ 10–34 J sec.
Solution Given Dt = 2.5 ¥ 10–14 sec.
Formula used is
 h
DE Dt = =
2 4p
h 1 6.62 ¥ 10 -34 1
DE = = ¥ = 0.211 ¥ 10 -20 J
4p Dt 4 ¥ 3.14 2.5 ¥ 10 -14
DE = 2.11 ¥ 10 -21 J
ExamplE 12 An excited atom has an average life-time of 10–8 sec. During this time period it emits a photon and
returns to the ground state. What is the minimum uncertainty in the frequency of this photon?
Solution Given Dt = 10–8 sec.
Formula used is
 h
DE Dt = =
2 4p
As E = hn or DE = D (hn ) = hDn
h 1
or hDvDt = or DvDt =
4p 4p
1 1 1 1
or Dv = = ¥
4p Dt 4 ¥ 3.14 10 -8
Dv = 7.96 ¥ 106 sec
E xamplE 13 Compare the uncertainties in velocity of a proton and an electron contained in a 20 Å box.
Solution Given Dx = 2.0 ¥ 10–9 m.

Formula used is
 h h
DpDx = = or Dp =
2 4p 4pDx
As uncertainty in momentum for electron and proton does not depend upon mass, we have
Dp
Dp = D (mv) = mDv or Dv =
m
As Dp p = Dpe
Dp p
Dv p =
mp
Dpe
and Dve =
me
Dv p Dp p me me
= = [ Dp p = Dpe ]
Dve Dpe m p m p
9.1 ¥ 10-31
= = 5.45 ¥ 10 -4
1.67 ¥ 10 -27
E xamplE 14 Find the energy of an electron moving in one dimension in an infinitely high potential box of
width 1.0 Å. Given m = 9.1 ¥ 10–31 kg and h = 6.62 ¥ 10–34 J sec.
Solution Given l = 1.0 ¥ 10–10 m, m = 9.1 ¥ 10–31 kg and h = 6.62 ¥ 10–34 J sec.
Formula used is
n2h2
En =
8mL2
n 2 (6.62 ¥ 10 -34 ) 2
=
8 ¥ 9.1 ¥ 10 -31 ¥ (1.0 ¥ 10 -10 ) 2
= 0.602 ¥ 10 -17 n 2 J
for n = 1,
E1 = 6.02 ¥ 10–18 J
and for n = 2,
E2 = 6.02 ¥ 10–18 ¥ 4 J
= 2.408 ¥ 10–17 J
= 2.41 ¥ 10–17 J
E xamplE 15 Calculate the energy difference between the ground state and the first excited state for an
electron in a box of length 1.0 Å.
Solution Given L =1.0 ¥ 10–10 m.
Formula used is
n2h2
En =
8mL2
Put n = 1 for ground state and n = 2 for first excited state
h2 (6.62 ¥ 10-34 ) 2 ¥ 3
E2 - E1 = [22 - 12 ] =
8mL 2
8 ¥ 9.1 ¥ 10-31 ¥ (1.0 ¥ 10-10 ) 2
= 1.81 ¥ 10-17 J
E xamplE 16 Compute the energy of the lowest three levels for an electron in a square well of width 3 Å.
Solution Given L = 3 ¥ 10–10 m.
Formula used is
n2h2
En =
8mL2
Put n = 1, 2, 3 for first three levels, then
h2 (6.62 ¥ 10 -34 )2
E1 = =
8mL2 8 ¥ 9.1 ¥ 10 -31 ¥ (3 ¥ 10 -10 ) 2
= 6.688 ¥ 10 -19 J
= 6.7 ¥ 10 -19 J
E2 = 4 E1 = 2.68 ¥ 10 -18 J and
E3 = 9 E1 = 6.03 ¥ 10 -18 J
E xamplE 17 An electron is bound in one-dimensional potential box which has a width 2.5 × 10–10 m.
Assuming the height of the box to be infinite, calculate the lowest two permitted energy values of the electron.
Solution Given L = 2.5 ¥ 10–10 m.
Formula used is
n2h2
En =
8mL2
For lowest two permitted energy values of electrons, put n = 1 and 2. Then
for n = 1,
h2 (6.62 ¥ 10 -34 ) 2
E1 = =
8mL 2
8 ¥ 9.1 ¥ 10 -31 ¥ (2.5 ¥ 10 -10 ) 2
= 9.63 ¥ 10 -19 J
for n = 2,
(2) 2 ¥ (6.62 ¥ 10 -34 ) 2
and E2 =
8 ¥ 9.1 ¥ 10 -31 ¥ (2.5 ¥ 10 -10 ) 2
= 3.853 ¥ 10 -18 J
E xamplE 18 Compute the lowest energy of a neutron confined to the nucleus which is considered as a box
with a size of 10–14 m.
Solution Given L = 10–4 m, h = 6.62 ¥ 10–34 J sec and m = 1.67 ¥ 10–27 kg.
Formula used is
n2h2
En =
8mL2
For lowest energy n = 1

h2 (6.62 ¥ 10 -34 ) 2
E1 = =
8mL2 8 ¥ 1.67 ¥ 10 -27 ¥ (10 -14 ) 2
E1 = 3.28 ¥ 10 -13 J
E xamplE 19 State the values of momentum and energy of a particle in one-dimensional box with impenetrable
walls. Find their values for an electron in a box of length 1.0 Å for n = 1 and n = 2 energy states. Given m =
9.1 ¥ 10–31 kg and h = 6.63 ¥ 10–34 J sec.
Solution Given L = 1.0 ¥ 10–10 m, n = 1 and 2, m = 9.1 ¥ 10–31 kg and h = 6.63 ¥ 10–34 J sec.
The Formulae used are
nh
pn = (i)
2L
pn2 n 2 h 2
En = = (ii)
2m 8mL2
Momentum for n = 1 and 2 are
1 ¥ 6.63 ¥ 10 -34
p1 = = 3.315 ¥ 10 -24 kg m/sec
2 ¥ 10 -10
2 ¥ 6.63 ¥ 10 -34
and p2 = = 6.63 ¥ 10 -24 kg m/sec
2 ¥ 10 -10
Energy for n = 1 and 2 are

n2h2 (1) 2 (6.63 ¥ 10-34 )2
E1 = =
8mL2
8 ¥ 9.1 ¥ 10 -31 ¥ (10 -10 )2
= 6.04 ¥ 10 -18 J
For n = 2,
n2h2
E2 = = (2) 2 E1 = 2.416 ¥ 10-17
8mL2
= 2.42 10-17 J
ExamplE 20 An electron is constrained to move in a one dimensional box of length 0.1 nm. Find the first three
energy eigen values and the corresponding deBroglie wave lengths. Given h = 6.63 ¥ 10–34 J sec.
Solution Given L = 1.0 ¥ 10–10 m.
Formulae used are
n2h2 nh
En = and pn =
8mL2 2L
for n = 1
h2 (6.63 ¥ 10 -34 )2
E1 = =
8mL2 8 ¥ (9.1 ¥ 10 -31 ) ¥ (10-10 )2
= 6.04 ¥ 10 -18 J
Similarly, for n = 2, is equal to

E2 = (2)2 E1 = 2.416 ¥ 10–17 J
and for n = 3,
E3 = (3)2 E1 = 5.36 ¥ 10–17 J
As we know
h nh
ln = and pn =
pn 2L
2L
or ln =
n
For n = 1
l1 = 2L = 2.0 ¥ 10–10 m = 2 Å
For n = 2,
2L
l2 = = l = 1.0 ¥ 10 -10 m = 1.0 Å
2
For n = 3,
2L
l3 = = 0.667 ¥ 10 -10 m = 0.667 Å
3
ExamplE 21 The minimum energy possible for a particle entrapped in a one dimensional box is 3.2 ¥ 10–18 J.
What are the next three energies in eV the particle can have?
Solution Given E1 = 3.2 ¥ 10–18 J.
Formula used is
n2h2
En = or En µ n 2 (i)
8mL2
32.0 ¥ 10 -19
Now energy in eV =
1.6 ¥ 10 -19
E1 = 20 eV
Next three values of energy can be obtained by putting n = 2, 3 and 4.

E2 = n2E1 = (2)2 E1 = 4 ¥ 20 eV = 80 eV
E3 = (3)2E1 = 9 ¥ E1 = 9 ¥ 20 eV = 180 eV
and E4 = (4)2E1 = 16 ¥ 20 eV = 320 eV
E xamplE 22 The energy of an electron constrained to move in a one dimensional box of length 4.0 Å is
9.664 ¥ 10–17 J. Find out the order of excited state and the momentum of the electron in that state. Given h =
6.63 ¥ 10–34 J sec.
Solution Given En = 9.664 ¥ 10–17 J and L = 4 ¥ 10–10 m
Formulae used are
n2h2 nh
En = and pn =
8mL2 2L
h2 (6.63 ¥ 10-34 )2
E1 = =
8mL 2
8 ¥ 9.1 ¥ 10-31 ¥ (4 ¥ 10-10 )2
= 3.774 ¥ 10 -19 J
n2h2
En = = n 2 E1
8mL2
E 966.4 ¥ 10 -19 J
or n2 = n =
E1 3.774 ¥ 10-19 J
or n = 16 (order of excited state)
Momentum of electron for n = 16

nh 16 ¥ 6.63 ¥ 10 -34
pn = =
2L 2 ¥ 4 ¥ 10 -10
= 13.26 ¥ 10 -24 kg m/sec
E xamplE 23 Evaluate the first three energy levels of an electron enclosed in a box of width 10 Å. Compare it
with those of glass marble of mass 1.0 gm, contained in a box of width 20 cm. Can these levels of the marble
be measured experimentally?
Solution Given for an electron n = 1 and L = 1.0 ¥ 10–9 m and for glass marble n = 1, L = 0.2 m and m = 1.0 ¥ 10–3 kg.
n2h2
Formula used is En =
8mL2
For electron
h2 (6.62 ¥ 10 -34 ) 2
E1 = =
8mL 2
8 ¥ 9.1 ¥ 10 -31 ¥ (1.0 ¥ 10 -9 ) 2
= 6.02 ¥ 10 -20 J
Similarly
E2 = (2)2E1 = 4 ¥ E1 = 24.08 ¥ 10–20 J
and E3 = (3)2E1 = 9 ¥ 6.02 ¥ 10–20 J
and = 54.18 ¥ 10–20 J
For glass marble
(6.62 ¥ 10 -34 ) 2
E1 =
8 ¥ 10 -3 ¥ (0.2) 2
= 1.3695 ¥ 10 -63 J
= 1.37 ¥ 10 -63 J
Similarly,
E2 = (2)2E1 = 5.48 ¥ 10–63 J
and E3 = (3)2E1 = 9E1 = 12.33 ¥ 10–63 J
It is clear that the levels in case of marble are very small and are nearly zero. So it is not possible to measure them
experimentally.
E xamplE 24 Find the smallest possible uncertainty in position of the electron moving the velocity 3 ¥ 107
m/sec (Given, h = 6.63 ¥ 10–34 J.sec, m0 = 9.1 ¥ 10–31 kg)
Solution By using the formula,
h
DXmin Dpmax =
2p
m0v
Dpmax = p = mv =
v2
1-
c
m0v h
Dxmin ¥ =
v 2 2p
1-
c2
2
v2 Ê 3 ¥ 107 ˆ
1- 2 6.63 ¥ 10-34 1 - Á ˜
c = Ë 3 ¥ 108 ¯
Dxmin = h
2p m0v 2 ¥ 3.14 ¥ 9.1 ¥ 10-31 ¥ 3 ¥ 107
12
= 3.8 10 m
E xamplE 25 Show that the uncertainty in the location of the particle is equal to deBrogile wave length the
uncertainty in its velocity is equal to the velocity.
Solution Given, Dx = A
Formula used is
DxDpx = h or xDp = h
h h
Dpx = = = px [{ x = l (given)]
x l
Dpx = px
or mDvx = mvx
or Dvx = vx Hence proved.
E xamplE 26 An electron is confined to move between two right walls separated by 10–9 m. Find the deBrogile
wavelengths respresenting the first three allowed energy state of the electron and the corresponding energies
(electron mass is 9.1 ¥ 10–31 kg and h = 6.63 ¥ 10–34 J sec.
Solution Given, L = 10–9 m = 10 Å, me = 9.1 ¥ 10–31 kg and h = 6.63 ¥ 10–34 J ◊ sec.
The electron moves forth and back between rigid walls will form a stationary wave-pattern with node at the walls, for this,
the distance L between the wall must be a whole multiple of the debrogile half-wavelengths.
l
Thus, L= n , where n = 1, 2, 3, …
2
2 L 2 ¥ 10 Å
l= =
n n
l1 = 20 Å, l2 = 10 Å, l3 = 6.7 Å
The corresponding energies are given by
nL h 2 n 2 ¥ (6.63 ¥ 10-34 ) 2
En = =
8 mL 8 ¥ 9.1 ¥ 10-31 ¥ (10-9 ) 2
2
= n 2 ¥ 6.04 ¥ 10-20 Joule

n 2 ¥ 6.04 ¥ 10-20
= = 0.38 n 2eV
1.6 ¥ 10-19
For n = 1, 2, 3, …
E1 = 0.38 eV, E2 = 1.52 eV, E3 = 3.45 eV
p p
E xamplE 27 The wave function of a certain particle is y = A cos2 x for - < x <
(i) Find the value of A. 2 2
p
(ii) Find the probability that the particle be found between x = 0 and x =
4
Solution
p p
(i) Given, y = A cos2 x for - <x<
2 2
By using the condition for normalisation
p /2 p /2
Ú yy 1dx = 1 or 2 A2 Ú cos
4
x dx = 1
- p /2 0
Èp p ˘ 3p
2 A2 Í + ˙ = 1 fi 2 A2 ¥ =1
Î8 8˚ 16
8
A
3
(ii) Probability is given by
p /4 p /4
Ú |y | Ú
2
P= dx = A2 cos 4 x dx
0 0
p /4 p /4
8
= A2 Ú cos 4 x dx =
3p Ú cos 4 x dx
0 0
8 È 3p + 8 ˘ 3p + 8 È 1 2 ˘
= = =Í + ˙
3p ÍÎ 32 ˙˚ 12p Î 4 37 ˚
P = 0.25 + 0.2123 = 0.4623
E xamplE 28 Normalise the wave function

y(x) = 0 outside the box of size l
y(x) = A sin kx for 0 < x < l
p
where k=
l
p
Solution Given, y(x) = A sin kx and k =
l
The condition for normalisation for the confinement 0 to l
Ú |y |
2
dx = 1
0
l
ÚA
2
For given problem, sin 2 kx = 1
0
On solving the above eqs,
2
A=
l
2 px
Then, y ( x) = sin
l l
E xamplE 29 Calculate the energy difference between the ground state and the first excited state for an
electron in one-dimensional rigid box of length 10–8 cm. (Mass of electron is 9.1 ¥ 10–31 kg and h = 6.63 ¥
10–34 J ◊ sec).
Solution Given, me = 9.1 ¥ 10–31 kg, h = 6.63 ¥ 10–34 J sec and L = 10–8 cm.
The energy of particle of mass m in 1-D rigid box of side L is given by,
n 2 L2
En = , n = 1, 2, 3, º
8 mL2
n 2 (6.63 ¥ 10-34 )2
= = 6 ¥ 10-18 ¥ n 2 Joule
8 ¥ 9.1 ¥ 10-31 ¥ (10-10 )2
6 ¥ 10-18 2
= n eV
1.6 ¥ 10-19
For ground state, n=1
E1 = 38 eV
For first excited state, n = 2
E2 = 152 eV
The energy difference DE = E2 – E1
= 152 – 38
= 114 eV
Q.1 Which of the following relations is correct for Heisenberg’s uncertainty principle?
 h 
(a) DE Dt ≥ (b) Dx Dp ≥ (c) DL Dq ≥ (d) All of these
2 4p 2
Q.2 Heisenberg uncertainty relation holds good for
(a) microscopic as well as macroscopic particles both
(b) only microscopic particles (c) only macroscopic particles
(d) none of these
Q.3 The energy of a particle in infinite potential well is
(a) proportional to n2 (b) inversely proportional to n2
(c) proportional to n (d) inversely proportional to n
Q.4 The momentum of a particle in infinite potential well of length l is

(a) proportional to l (b) inversely proportional to l
(c) proportional to l2 (d) inversely proportional to l2
Q.5 The momentum of a particle in infinite potential well is
(a) proportional to n (b) inversely proportional to n2
2
(c) proportional to n (d) inversely proportional to n
Q.6 Which one of the following operators is associated with energy
2 2  2 2  ∂
(a) - — +V (b) - — (c) — (d) i
2m 2m i ∂t
Q.7 Which one of the following operators is associated with kinetic energy
2 2 2 2 ∂  2
(a) - — (b) - — +V (c) i (d) - —
2m 2m ∂x 2m
Q.8 Which one of the following operators is associated with momentum
 2 2 2 2 2 2
(a) — (b) - — (c) - — (d) - —
i i i 2m
Q.9 Which one of the following energy values of a particle in infinite potential well of length l is allowed
2 2 2 n 2p 2  2 p 2 2 n2 2
(a) n p  (b) (c) 2 (d)
2ml 2ml 2 2mn l 2ml 2
Q.10 Which one of the following is the radius of first Bohr’s orbit
e 0 h2 e 0 2 e0 e0
(a) 2 (b) 2 (c) 2 (d)
p mZe p mZ e p mZe p m2 Z 2 e2
2
Q.11 The entire information of a quantum system can be gathered with the help of
(a) position (b) eigen value
(c) momentum operator (d) wave function
Q.12 The expression |y(x, t)|2 stands for
(a) normalisation (b) position
(c) time probability density (d) position probability density
+•
Q.13 If y is normalised wave function, then the value of Ú– • y * y dV will be
(a) zero (b) 1 (c) • (d) – •
Q.14 Which of the following relation is correct for Schrödinger’s wave equation, moving also x-axis?
∂2y 2m ∂2y 2m
(a) + 2 ( E - V )y = 0 (b) - 2 ( E - V )y = 0
∂x 2  ∂x 2 
∂2y 2m 2
(c) + 2 ( E - V )y = 0 (d) none of these
∂x 2 
Q.15 The wave function ‘y’ associated with matter waves has no direct physical significance. It
(a) is a complex quantity (b) is not an observable quantity
Q.16 The normalised eigen wave function of a particle in a box of length ‘L’ is
2 np x 2 np x 2 np x
(a) sin (b) sin (c) sin (d) none of these
L L L L L L
Q.17 The energy levels of a particle in a box are
(a) equally spaced (b) continuous
(c) not-equally spaced (d) none of these
Q.1 What is Heisenberg’s uncertainty principle?

Q.2 What do you understand by wave function?
Q.3 What is normalised wave function?
Q.4 Write and explain Schrödinger’s time dependent and time independent equations.
Q.5 What are the basic assumptions of Bose-Einstein statistics?
Q.6 Differentiate between classical and Quantum statistics.
Q.7 What do you understand by normalised and orthogonal wave functions?
P RACtiCe P RoBleMs
general questions
Q.1 Starting from deBroglie’s wave concept obtain Heisenberg’s uncertainty principle.
Q.2 State Heisenberg’s uncertainty principle and derive it from a hypothetical gamma ray microscope.
Q.3 Why is uncertainty principle important for microscopic particles but significant in practical life?
Q.4 Illustrate Heisenberg’s uncertainty principle by diffraction of a beam of electrons by a narrow slit.
OR
Prove position momentum uncertainty principle using particle approach.
Q.5 By applying uncertainty principle explain non-existence of electrons in atomic nucleus.
Q.6 What other reasons show why electrons cannot exist inside the nucleus?
Q.7 Apply Hiesenberg’s uncertainty principle to explain the following.
(a) Non-existence of electrons within the nucleus
(b) Existence of protons, neutrons and a-particles
(c) Existence of finite zero-point energy
(d) Binding energy of an electron in a hydrogen atom is of the order of 15 eV
Q.8 Explain the difference between quantum mechanics and classical mechanics.
Q.9 What do you understand by the wave function y of a moving particle?
Q.10 Give the physical significance of wave function. What does the square of wave function signify?
Q.11 What are the conditions and limitations that the wave function must obey?
Q.12 Starting from the wave equation and introducing energy and momentum of the particle obtain an
expression for three dimensional Schrödinger’s equation in time dependent form.
Q.13 Obtain three dimensional time independent Schrödinger’s wave equation from time dependent
Schrödinger’s equation.
Q.14 Derive an expression for Schrödinger time independent and time dependent wave equations.
Q.15 Derive time dependent Schrödinger wave equation.
Q.16 Give the formulation of time dependent Schrödinger equation for a free particle. Discuss the
interpretation of position, probability density and normalisation of wave function.
Q.17 Derive both time independent and time dependent Schrödinger equations for non-relativistic particle.
dy
Q.18 Why should y and be continuous everywhere?
dx
Q.19 What do you understand by orthogonal wave function? Explain orthogonality and orthonormality of
wave functions.
Q.20 Obtain Schrödinger’s wave equation for a particle in square well potential and discuss energy levels
when the well is infinitely deep.
Q.21 Discuss quantum mechanically the problem of linear harmonic oscillator and obtain its eigen values.
Also, write significance of zero point energy.
Q.22 Why are we not aware of quantisation in daily experience? Explain it.
Q.1 An electron of mass 9.1 × 10–31 kg has a speed of 1.0 m/sec with an accuracy of 0.05%. Calculate the
uncertainty with which the position of the electron can be located. [Ans: 1.15 ¥ 10–4 m]
Q.2 The electron in hydrogen atom may be confined to a nucleus of radius 5 × 10–11 m. Find out the
minimum uncertainty in the momentum of the electron and also find out the minimum kinetic energy
of the electron. Given m = 9.0 × 10–31 kg and h = 6.62 × 10–34 J sec.
[Ans: 1.054 ¥ 10–24 kg m/sec, 6.142 J]
Q.3 The speed of a bullet of mass 50 gm is measured to be 300 m/sec with an uncertainty of 0.01%. With
what accuracy can we locate the position of the bullet if it is measured simultaneously with its speed.
[Ans: 3.5 ¥ 10–32 m]
Q.4 Life time of a nucleus in the excited state is 10–12 sec. Calculate the probable uncertainty in energy and
frequency of a g-ray photon emitted by it. [Ans: 1.054 ¥ 10–22 J; 1.59 ¥ 1011 Hz]
Q.5 Compute the energy difference between the ground state and first excited state for an electron in a one-
dimensional rigid box of length 10–8 cm. Given m = 9.1 × 10–31 kg and h = 6.626 × 10–34 J sec
[Ans: 114 eV]
Q.6 Calculate the value of lowest energy of an electron in one dimensional force free region of length 4 Å.
[Ans: 3.78 ¥ 10–19 J]
Q.7 The lowest energy possible for a certain particle entrapped in a box is 40 eV. What are the next three
higher energies the particle can have? [Ans: 160 eV, 360 eV and 640 eV]
Q.8 Find the energy levels of an electron in a box 1 nm wide. Mass of electron is 9.1 × 10–31 kg. Also find
the energy levels of 10 gm marble in a box 10 cm wide.
[Ans: 6.02 ¥ 10–20 J, 24.08 ¥ 10–20 J and 54.18 ¥ 10–20 J; and for marble 5.49 ¥ 10–64 J,
21.96 ¥ 10–64 J and 49.41 ¥ 10–64 J]
Free Electron Theory 17

Learning Objectives
LO 1 Understanding Lorentz-Drude theory LO 3 Explain quantum theory of free electrons
LO 2 Learn about application and limitation LO 4 Discuss thermionic emission,
of Lorentz-Drude theory Richardson’s equation
Introduction
The simplest metals are alkali metals, which include sodium (Na), potassium (K), lithium (Li) etc. The
electronic configuration of Na atom is 1s2, 2s2, 2p6, 3s1. Thus, the valance electron is in the 3s state. This
electron behaves as conduction electron in the metal. The remaining 10 electrons of Na+ ion core fill the
1s, 2s and 2p states, which contain 2, 2 and 6 electrons, respectively. The distribution of core electrons is
the same as in the free ion in the metal. This way we can say that the metal crystal contains the positive
ion with the free electrons. These free electrons behave like the molecules in a perfect gas and are called
free electron gas.
According to free electron theory, a metal can be considered to consist of ion cores having the nucleus
and electrons other than valence electrons. These valence electrons form an electron gas, surround
the ion cores and are free to move anywhere within the metal. Thus, the valence electrons of the atom
become conduction electrons. In the theory, the force between the conduction electrons and ion cores
are neglected so that the total energy of the electron is all kinetic, i.e., the potential energy is taken to be
zero. Hence, the motion of the electrons within the metal is free because there are no collisions, similar
to the molecules of an ideal gas.
Long back, it was believed that many physical properties of metals including electrical and thermal
conductivities can be understood by considering free electron model. Attempts have been made by Drude
and Lorentz to explain quantitatively the conductivities of metals on the basis of free electron theory.
In this context it is necessary to understand first the main characteristics of metals, which are discussed
below:

(i) Metals obey Ohm’s law, i.e., in steady state, the current density J is proportional to the applied

electric field strength E . It means
   
J µ E or J = s E
where, s is electrical conductivity.

(ii) Metals possess high electrical and thermal conductivities.
(iii) At low temperature, the resistivity r is proportional to the fifth power of absolute temperature, i.e.,
1
r µ T 5 , where r =
s
(iv) The resistance of metals increases with rise in temperature, i.e., they have positive temperature
coefficient. The resistance of certain metals vanishes at absolute zero and they exhibit the
phenomenon of superconductivity.
1
(v) For most of the metals, resistivity is inversely proportional to the pressure, i.e., r µ .
P
(vi) The ratio of thermal and electrical conductivities is directly proportional to the absolute
temperature and this phenomenon is known as Wiedemann–Franz Law, i.e.,
K K
µ T or = constant
s sT
17.1 Lorentz–DruDe theory: CLassiCaL Free eLeCtron LO1

theory oF MetaLs
As mentioned earlier, the valence electrons in a metal form an electron gas which surrounds the ion cores and
are free to move throughout the crystal lattice and behave as molecules in the perfect gas. These electrons
are negatively charged. Therefore, the flow of electricity corresponds to their motion. But in the absence of
externally applied field the electrons are moving in the random directions and collide with the residual ions
frequently. Thus, the laws of classical theory of gases can be applied to free electron gas also. This way, the
electrons can be assigned a mean free path l, mean collision time t and average velocity v–.
In 1900, Drude used the electron-gas model to explain electrical conductivity in metals theoretically. Later Drude
theory was modified by Lorentz and named as Lorentz-Drude free-electron theory. Lorentz and Drude assumed
that at ordinary temperature, free electrons move in metals randomly with an average speed of the order of
105 m/sec. During the random motion, these electrons collide with themselves and atoms or ions of lattice and
have no practical contribution to electrical and thermal conductivities Hence, in the absence of external electric
field, the contribution of electrons to current in metal will be zero. If an external electric field is applied to the
metals, the electrons will be accelerated in the opposite direction to the applied electric field and produce a
current. In thermal equilibrium, the free electrons are assumed to follow Maxwell–Boltzmann distribution.
17.2 appLiCations oF Lorentz–DruDe theory LO2

The Lorentz–Drude theory can be applied to explain the properties of metals including their electrical and
thermal conductivities.
17.2.1 electrical Conductivity
We consider that there are n free electrons per cubic meter in the metal. If we apply an electric field to the
metal, the electrons modify their random motion and move with an average drift velocity vd in the opposite
direction to that of the applied field. The magnitude of force experienced by the electron is given by
eE = ma (i)
eE
or a=
m
Thus, the electrons undergo an acceleration eE/m. However, the electron will not accelerate indefinitely but
after a short period it collides with a +ve ion in the metal. At each collision its velocity is reduced to zero. So
it is accelerated between two collisions only. If l is the mean free path and t is the free time, then the time
taken between two successive collisions
l
t= (ii)
v
eE
During this time the velocity will be t . Thus the velocity at the beginning of the path is zero and at its end
m
eE
is t . Hence, the average drift velocity vd will be the mean of the two, i.e.,
m
1 eE
vd = t (iii)
2 m
The current density is then given by
È eE ˘
J = nevd = ne Í t
Î 2m ˙˚
ne 2t
or J= E (iv)
2m
Metals obey Ohm’s law which states that in steady state the current density J is proportional to the electric
field strength.
JµE or J = sE (v)
By using Eqs. (iv) and (v), we get
ne 2t (vi)
s=
2m
By putting the value of t from Eq. (ii) in Eq. (vi) we get
ne 2 l
s= (vii)
2mv
Since the energy

1 2 3
mv = kT
2 2
ne 2 l v
s=
6kT
17.2.2 thermal Conductivity
Free electrons also contribute to the conduction of heat energy in the Y
metals. As we have already understood, free electrons behave similar to the
molecules of a perfect gas. They possess greater kinetic energy at the hot P1
end of the metal sheet than at the cold end. Suppose n number of electrons
are moving randomly in all directions in the metal and their motion can be P
n
resolved along the three axes. Along any one particular direction only X
6 P2
electrons will move, as there are six directions of their possible motions
along three axes. We consider three planes of unit area such that the plane
P1 and P2 are at the same distance of mean free path l from the plane P.
Z
The temperatures of planes P1 and P2 are T1 and T2, respectively. If the
Figure 17.1
temperature of plate P1 is greater than that of plate P2, the energy will
nv
transfer from plate P1 to P2 i.e., electrons will transfer from P1 to P2. Since each electron has energy
1 2 3 6
mn = kT1, (k is the Boltzmann constant), the net energy transfer from P1 to P2 per unit area per second
2 2
will be
nv 3
¥ kT1 (i)
6 2
and the same way, the net energy transfer from P2 to P1 will be
nv 3
¥ kT2 (ii)
6 2
Thus, the net transfer of energy from plate P1 to plate P2 through plate P per unit area per second will be
nv È 3kT1 3kT2 ˘
Q= -
6 ÍÎ 2 2 ˙˚
nvk
= [T1 - T2 ]
4 (iii)
[T1 - T2 ]
But Q= K (iv)
2l
where K is thermal conductivity of metal.
By using Eqs. (iii) and (iv), we have

[T1 - T2 ]K nvk
= [T1 - T2 ]
2l 4
knvl (v)
or K=
2
The above equation represents the expression for thermal conductivity of a metal.
17.2.3 Wiedemann-Franz Law
As we have deduced the expressions for thermal conductivity K and the electrical conductivity s, now we are
in the position to prove Wiedemann-Franz law which states that ratio of thermal conductivity K to the electrical
conductivity s is proportional to the absolute temperature. From the expressions of K and s, we have
K knvl /2 Ê ne 2 l ne 2 l v ˆ
= 2 ÁË s = = ˜
s ne vl /6kT 2mv 6kT ¯
2
K Ê kˆ (i)
= 3Á ˜ T
s Ë e¯
2
or
K Ê kˆ (ii)
= 3Á ˜
sT Ë e¯
Putting the values of Boltzmann constant k and the charge of electron, we find
2
K Ê 1.38 ¥ 10-23 ˆ -8
= 3Á ˜ = 2.23 ¥ 10
sT Ë 1.6 ¥ 10-19 ¯
K
Thus, has the same values at all the temperatures for all the metals or the ratio K/s is directly proportional
sT
to absolute temperature. This is called Weidemann-Franz Law.
17.3 LiMitations oF Lorentz–DruDe or Free eLeCtron theory LO2

The limitations of free electron theory are listed below.
(i) It does not explain why only some crystals are metallic.
(ii) It does not explain why the metals prefer only certain structures.
(iii) In real situations, the electrical conductivity depends on the temperature. The free electron theory
does not explain the temperature variation of electrical conductivity.
(iv) The paramagnetism of metals is nearly independent of temperature. This result could also not be
explained by this theory.
17.4 QuantuM theory oF Free eLeCtrons LO3

As we know that the metal contains a large number of conduction electrons which are not completely free (but
partially), though they are not bound to any particular atomic system. The forces between conduction electrons
and ion cores are neglected in the free electron approximation so that the electrons within the metal are treated
as free. Further, the energy possessed by electron is kinetic, since the potential energy is taken to be zero.
Consider an electron of mass m confined in a box of length L. Under this situation, the Schröedinger wave
equation becomes
2 2
- — y = Ey (i)
2m
The solution of the above equation is
 
y = y 0 exp(i k ◊ r ) (ii)
 2p
where k is the wave vector with the magnitude k = .
l
It can be shown from Eq. (ii) that
∂2y ∂2y ∂2y
2
= - k x2y ; 2
= - k y2y ; 2
= - k z2y
∂x ∂y ∂z
Then,
∂2y ∂2y ∂2y
— 2y = + +
∂x 2 ∂y 2 ∂z 2
= - ( k x2 + k y2 + k z2 )y (iii)
or =2y = –k2y (iv)

By using Eqs. (i) and (iv), we have
2
- ¥ ( - k 2y ) = Ey
2m
2 k 2 (v)
or E(k ) =
2m
Now we can write the total energy E in terms of momentum p with the help of following relations.
Since, 2p h
k= and l =
l p
2 k 2 h 2 1 4p 2
E= =
2m 4p 2 2m l 2
h2 1 h2 p 2 p2
= 2
= 2
=
2m l 2m h 2m
p2
or E= (vi)
2m
Eq. (vi) represents the energy of a free particle (i.e., electron) and thus the energy is continuous. Here it may be
mentioned that we have not considered the lattice periodicity and also assumed the constant potential inside the
2p n
crystal to be zero. However, for cyclic boundary conditions, k = , where L is the length of the cyclic chain
L
(i.e., the solid). Therefore
2 k 2 L = 3l/2
E (n) =
Quantum Number n →
2m
n2 h2 (vii)
E (n) =
2mL2 L = 2l/2
The first three lower energy state wavefunctions are represented in L = l/2
Fig. 17.2. The distribution of the available electrons among the various
allowed energy levels and the evaluation of the related quantities can
be understood better along with the treatment of the free electron gas in
three-dimensional box of length L.
x=0 x→ x=L
17.4.1 Fermi energy Figure 17.2
Consider that N free electrons are contained in a box at absolute temperature. At ky
absolute zero all the energy levels below a certain level will be filled with electrons and
the levels above this level will be empty. The energy level which divides the filled and
empty levels is called ‘Fermi level’ and the corresponding energy of that level is known kF
as ‘Fermi-energy’ EF. In ground state of the system of N free electrons, the occupied
states may be represented as a point inside a sphere in k-space as shown in Fig. 17.3. kx
The kx, ky and kz are the components of kF along X, Y and Z axes, respectively. As per
previous article, the energy of the electron is given by
2 k 2 kz
Ek = (i)
2m Figure 17.3
From the above relation it is clear that the energy increases as the square of distance from the origin of the k
space coordinate system. All the electrons which lie on the same spherical shell of radius, kF, have the same
energy, which is called Fermi Energy. It is given by
2 2
EF = kF (ii)
2m
2p
Since, k= n ky
L
2p 4p/L
kx = nx
L
2p 2p/L
ky = ny
L
2p
kz = nz kx
L 2p/L 4p/L 6p/L
where nx, ny and nz have the values 0, ±1, ±2, … . Therefore
2p 4p 6p 2p/L
k x = 0, ± ,± ,± ,º
L L L
3
ky and kz also have the same values. Suppose ÊÁ ˆ˜ is the
2p 4p/L
Ë L¯ kz
volume of one shell in k-space (Fig. 17.4). Then in a sphere of Figure 17.4
volume 4p k F3 , the total number of energy states (or shells) will be

3
4p 3
kF
NF = 2 3 3 (iii)
Ê 2p ˆ
ÁË ˜¯
L
1 1
Here the factor 2 represents the two allowed values of spin quantum number ms = + (spin up) and ms = -
2 2
(spin down) for each allowed value of k.
If all of these energy shells are filled with electrons, then N will be equal to the number of electrons, i.e.,
4p 3
kF
V
2 3 3 =N or N= k F3 ( V = L3) (iv)
Ê 2p ˆ 3p 2
ÁË ˜¯
L
Ê 3N p 2 ˆ
1/ 3 (v)
3N p 2
or k F3 = or k F = Á ˜
V Ë V ¯
From Eq. (v), it is clear that kF depends upon electron concentration ÊÁ ˆ˜ or in other wards kF depends upon
N
ËV ¯
number of electrons per unit volume but it does not depend on the mass of electrons. Now the Fermi energy is
2 2
EF = kF
2m
The energy can be written as
1 2
EF = mvF (vii)
2
where vF is the velocity of electron in Fermi level, i.e., corresponding to Fermi energy. Then
2/ 3
1 2 2 È 3N p 2 ˘
mvF = EF = Í ˙
2 2m Î V ˚
1/ 3
 È 3N p 2 ˘
\ vF = Í ˙ (viii)
mÎ V ˚
17.4.2 effect of temperature on Fermi–Dirac Distribution

According to Fermi-Dirac distribution law, the most probable distribution is given by
g(E)
n( E ) = (a + b E ) (i)
e +1
- EF 1
where a = and b = together with k as Boltzmann constant. It is often convenient to introduce the
kT kT
Fermi-distribution function f(E), which is defined as
n( E ) 1
f (E) = =
g ( E ) e(a + b E ) + 1
With the values of a and b, this function can be written together with EF as the Fermi energy
1
f ( E ) = ( E - E ) / kT
e F +1
At absolute zero (T = 0)
E - EF
= – •, if E < EF
kT
= +• if E > EF
and the Fermi distribution function
1
f (E) = –•
=1 for E < EF
e +1
1
= =0 for E > EF
e +• + 1
At any temperature T, and for E = EF,

1 1 1
f (E) = ( E - EF ) kT
= =
e +1 0
e +1 1+1
1
or f (E) =
2
The energy level corresponding to E = EF is called Fermi level. Since at E = EF, f(E) = 1/2, the Fermi level
is defined as the energy level at which there is a 1/2 probability of finding an electron. It depends on the
distribution of energy levels and the number of electrons available.
17.4.3 Density of states
It is defined as the number of energy states per unit energy range. It is denoted by the symbol D(E). In other words,
the ‘density of states’ for electrons in a band gives the number of orbitals (or states) in a certain energy range.
Hence, the number of filled (i.e., density of electron state) states having the energy in the range E and E + dE is
N(E)dE = E(E) f (E)dE (i)
From the energy relation
2 k 2
E=
2m (ii)
V
and N= k3
3p 2
2/ 3
 2 È 3p 2 N ˘
we get E= Í ˙ (iii)
2m Î V ˚
where E is total energy and N is the number of electrons. Therefore, from Eq. (iii), we have
3/ 2 3/ 2
3p 2 N È 2mE ˘ È 2 mE ˘
=Í 2 ˙ =Í 2 2˙
V Î  ˚ Î h /4p ˚
8pV
or N= (2mE )3/ 2 (iv)
3h3
By differentiating Eq. (iv) w.r.t. E, we get
dN 8pV 3 8p mV (2mE )1/2
= 3 (2m )3/ 2 E1/ 2 =
dE 3h 2 h3
3/ 2
dN V È 2m ˘
or = ( E )1/ 2 (v)
dE 2p 2 ÍÎ  2 ˙˚
The quantity ÊÁ
dN ˆ
is frequency referred to as the density of available
Ë dE ˜¯ f(E)
state D(E), which on multiplication with probability of occupation f (E)
N(E)
gives density of occupied state N(E), as shown in Fig. 17.5.
Thus, the number of electrons whose energies lie between E and E + dE
is given by
dN
N ( E )dE = f ( E )dE
dE
E EF
8p mV dE
N ( E )dE = 3
(2mE )1/ 2 ( E - E ) / kT Figure 17.5
h e F +1
17.4.4 average Kinetic energy of Free electron Gas at 0 K

All the electrons have energy less than the Fermi energy EF at 0 K, i.e., E < EF. With this condition, the Fermi-
Dirac distribution function becomes
1
f (E) = ( E - EF ) / kT
=1
e +1
If we consider that the average energy of an electron is Ee . Since f (E)= 1, we can write
1 EF 1 EF dN
Ee =
N Ú0 EN ( E )dE =
N Ú0 E
dE
f ( E )dE
dN
Substituting the value of from Eq. (v) and f (E) as 1, in the above equation, we get
dE
1 Ê 8p mV ˆ 1/ 2 EF 3/ 2
Ee = Á ˜ (2m ) Ú0 E dE
N Ë h3 ¯
1 Ê 8p mV ˆ 2
Ee = ÁË 3 ˜
(2m )1/ 2 EF5/ 2
N h ¯ 5
Now the above relation with the value of N substituted from Eq. (iv) reads
3
Ee = EF
5
17.5 therMioniC eMission LO4

When a metal is heated, electrons are emitted from its surface. Y
This phenomenon is known as thermionic emission. It is well
f =(W–EF)
known that the electrons in a metal are not completely free, but
they are bound to the metal surface by some attractive forces. w
The term EF is the Fermi energy, i.e., the energy of electron in E
F
metal at absolute zero, which is insufficient to escape out an
electron from the metal surface (Fig. 17.6). According to free
electron theory, W (potential energy) is the minimum energy X
which is required to emit an electron from the metal surface.
Therefore, the minimum energy to be supplied to the electron
Z
for emission is (W – EF) which is known as work function.
Figure 17.6
Hence, the emission of electrons from the surface of the metal
is possible only if the energy is more than that of work function.
17.5.1 phase space
To understand phase space, first of all we should know about position space and momentum space. The
three dimensional space in which the location of a particle is completely specified by the three position
coordinates (x, y, z) is known as position space. The instantaneous motion of a particle is described by
velocity components vx, vy and vz. However, for many purposes it is more convenient to use the corresponding
momentum components px, py and pz. The three dimensional space in which the momentum of a particle is
determined by the three momentum coordinates is known as momentum space. The combination of position
space and momentum space is known as phase space. The phase space is a six dimensional space (x, y, z,
px, py, pz). If dv and dp are the elements of volume enclosed by any particular cell in position space and
momentum space, respectively, then
dt = dv dp
The elementary volume enclosed by this cell in the phase space is given by
dt = dx dy dz dpx dpy dpz
= dx dpx dy dpy dz dpz = h3 ( dx dpx = dy dpy = dz dpz = h)
Total volume in phase space V Ú Ú Ú dpx dp y ◊ dp y

Total number of cells in phase = =
Volume of one cell h3
Ú 4p p dp
2
= (i)
h3
17.5.2 richardson’s equation
Richardson’s equation is well-known equation used for thermionic emission. This equation enables us to find
the emission current density of electrons.
If we represent f = W – EF, where W is the work function and EF is the Fermi energy, then the emission
current density J is given by
J = AT 2e–f/kT (i)
where A is a constant given by
4p mek 2
A= , and T is the temperature of the metal.
h3
Equation (i) is known as Richardson’s equation. However, there have been various theoretical expressions for
the constant A based on different physical assumptions. A significant work was done by Dushman, Fowler,
Sommerfeld and Nordheim in addition to the remarkable work of Richardson. A modern theoretical treatment
given by Modinos assumes the band theory of the emitting material. Then according to the form of the
constant A, Eq. (i) is also known as Dushman’s equation, Richardson-Dushman equation, and Richardson-
Law-Dushman equation. The main understanding of this equation is that due to the exponential function the
current increases rapidly with temperature when kT is less than f. However, for essentially every material
melting occurs well before kT = f.
s UmmarY
✦ According to free electron theory, a metal can be considered to consist of ion cores having the nucleus
and electrons other than valence electrons. These valence electrons form an electron gas, surround
the ion cores and are free to move anywhere within the metal. Thus, the valence electrons of the atom
become conduction electrons. In the theory, the force between the conduction electrons and ion cores
are neglected so that the total energy of the electron is all kinetic, i.e., the potential energy is taken to be
zero. Hence, the motion of the electrons within the metal is free because there are no collisions, similar
to the molecules of an ideal gas.
✦ In the theory given by Lorentz-Drude, they assumed free electrons to move in metals randomly with an
average speed. During this random motion the electrons were considered to collide with themselves and
with atoms or ions of lattice. Further, it was assumed that these electrons have no practical contribution
to the electrical and thermal conductivities. However, in the presence of external electric field, these
electrons are accelerated and hence produce the current. In thermal equilibrium, these free electrons are
assumed to follow the Maxwell-Boltzmann distribution.
✦ Based on Lorentz-Drude theory, the electrical and thermal conductivities of the metals were explained.
✦ Weidemann-Franz law states that the ratio of thermal conductivity K to the electrical conductivity of
the metal is proportional to the absolute temperature.
✦ Limitations of free electron theory were talked about.
✦ In free electron approximation, the electrons within the metal are treated as free as the forces between
conduction electrons and ion cores are neglected. However, the free electron was treated as confined in
a box of length L as per quantum theory of free electrons. Using the Schröedinger equation, the energy
of the electron was calculated and it was realised that the energy is not continuous rather it is quantized.
Finally, the lower energy state wave functions were represented. However, the distribution of the
available electrons among the various allowed energy levels and the evaluation of the related quantities
can be understood better along with the treatment of the free electron gas in three-dimensional box.
✦ At absolute temperature, all the energy levels below a certain level are filled with electrons and the
levels above this level are empty. The energy level which demarcates the filled and empty levels is
called Fermi level. The energy corresponding to the Fermi level is referred to as Fermi energy.
✦ The concept of k space was given and the total number of energy sates (or shells) was calculated.
Finally, the velocity of electrons in Fermi level, i.e., corresponding to Fermi energy, was obtained as
1/ 3
 È 3N p 2 ˘
vF = Í ˙ , where N is the number of electrons and V the volume of the box of length L.
mÎ V ˚
✦ Effect of temperature on Fermi-Dirac distribution was discussed and then the concept of density of
states D(E) was given. The density of states is nothing but is the number of energy states per unit energy
range. In other words, the density of sates for electrons in a band gives the number of orbitals (or states)
in a certain energy range.
✦ Concept of phase space was introduced in order to explain the thermionic emission. The three-
dimensional space in which the momentum of a particle is determined by the three momentum
coordinates (px, py, pz) is called the momentum space. The combination of position space (x,y,z) and
the momentum space (px, py, pz) is known as the phase space. So the phase space is represented by (x,
y, z, px, py, pz).
✦ The Richardsons’s equation that enables us to find the emission current density of the electrons
was mentioned. It is represented as J = AT2e–f/kT, where f = W – EF together with W as
È 4p mek 2 ˘
the work function and EF as the Fermi energy. The constant A = Í ˙ together with k as the
Î h3 ˚
Boltzmann constant and m and e as the mass and charge of the electron, respectively.
s olVeD e XamPles
E xamplE 1 Determine the average energy and speed of electron at its mean energy at 0 K, if the Fermi
energy is 10 eV.
Solution Given EF = 10 eV
3
Average energy E0 = EF
5
3
= ¥ 10 eV
5
= 6.0 eV
1 2 2 E0
and, mv = E0 or v =
2 m
2 ¥ 6.0 ¥ 1.6 ¥ 10-19
v=
9.1 ¥ 10-31
v 1.45 ¥ 106 m / sec
ExamplE 2 Fermi energy of a given substance is 7.9 eV. What is the average energy and speed of electron in this
substance at 0 K?
Solution Given EF = 7.9 eV
3
Average energy E0 = EF
5
3
E0 = ¥ 7.9 eV
5
= 4.74 eV
1 2 2 E0
and, mv = E0 or v =
2 m
2 ¥ 4.74 ¥ 1.6 ¥ 10-19
or v= = 1.29 ¥ 106 m / sec
9.1 ¥ 10-31
E xamplE 3 There are 2.5 ¥ 1028 free electrons per cubic meter of sodium. Calculate the Fermi energy and
Fermi velocity.
Solution Given N = 2.5 ¥ 1028

V
Formula used is
2/3
1 2 2 Ê 2 N ˆ
EF = mvF = Á 3p ◊ ˜¯
2 2m Ë V
2/3
h2 Ê 2 N ˆ
or EF = Á 3p ◊ ˜¯
8p 2 m Ë V
(6.62 ¥ 10-34 ) 2
= [3 ¥ (3.14) 2 ¥ 2.5 ¥ 1028 ]2 / 3
8 ¥ (3.14) 2 ¥ 9.1 ¥ 10-31
= 4.99 ¥ 10-19 J
EF = 5.0 ¥ 10-19 J
or EF = 3.12 eV
1/ 3
h È 2 N˘
and Fermi velocity vF = 3p ◊ ˙
2p m ÍÎ V˚
6.62 ¥ 10-34
= (3 ¥ (3.14) 2 ¥ 2.5 ¥ 1028)1/ 3
2 ¥ 3.14 ¥ 9.1 ¥ 10-31
= 1.05 106 m / sec
E xamplE 4 The density of copper is 8940 kg/m3 and atomic energy weight is 63.55. Determine the Fermi
energy of copper. Also obtain the average energy of free electrons of copper at 0 K.
Solution Given atomic weight = 63.55 kg and density of copper = 8940 kg/m3.
63.55 kg
Volume of 1 kg mole of copper, V =
8940 kg/ m 2
Number of atoms per kg atom = 6.02 ¥ 1026
or N 6.02 ¥ 1026 ¥ 8940

=
V 63.55
2/3
1 2 2 È 2 N ˘
Fermi energy EF = mvF = 3p
2 2m ÍÎ V ˙˚
2/3
h2 È 2 N ˘
= 3p
8p 2 m ÍÎ V ˙˚
2/3
(6.62 ¥ 10 -34)2 È 2 6.02 ¥ 1026 ¥ 8940 ˘
= Í 3 ¥ (3.14) ¥ ˙
8 ¥ (3.14) 2 ¥ 9.1 ¥ 10-31 Î 63.55 ˚
= 11.261 ¥ 10-19 J
EF = 7.038 eV
Average energy E0 = 3 EF = 3 ¥ 7.038 eV

5 5
= 4.22 eV
E xamplE 5 Consider silver in the metallic state with one free electron per atom. Calculate the Fermi energy.
Given that density of silver is 10.5 g/cm3 and atomic weight is 108.
108 g
Solution Volume of 1 g mole of silver, V = and number of atoms per g atom = 6.02 ¥ 1023.
10.5 g/cm3
N 6.2 ¥ 1023 ¥ 10.5
=
V 108
= 5.85 ¥ 1022 per cm3
= 5.85 ¥ 1028 per cm3
2/3
h2 È 2 N ˘
Fermi energy EF = 3p ◊ ˙
8p 2 m ÍÎ V˚
(6.62 ¥ 10-34 ) 2
= [3 ¥ (3.14) 2 ¥ 5.85 ¥ 1028 ]2 / 3
8 ¥ (3.14) 2 ¥ 9.1 ¥ 10-31
= 8.799 ¥ 10-19 J
= 5.499 eV
= 5.5 eV
ExamplE 6 Aluminium metal crystallises in f.c.c. structure. If each atom contributes single electron as free
electron and the lattice constant a is 4.0 Å, treating conduction electron as free electron Fermi gas, find (i)
Fermi energy (EF) and Fermi vector (kF) and (ii) total kinetic energy of free electron gas per unit volume at
0 K.
Solution In f.c.c. lattice number of electrons per unit cell will be (N) = 4 and volume of a unit cell is a3 = 64 ¥ 10–30 m3
N 4
and =
V 64 ¥ 10-30
= 6.25 ¥ 1028
2 2/3
Fermi energy EF = h È3p 2 ◊ N ˘
8p 2 m ÍÎ V ˙˚
(6.62 ¥ 10-34 )2
= [3 ¥ (3.14) 2 ¥ 6.25 ¥ 1028 ]2 / 3
8 ¥ (3.14) 2 ¥ 9.1 ¥ 10-31
= 9.2 ¥ 10-19 J
= 5.75 eV
1/ 3
È 2 N˘
Fermi vector k F = Í3p ◊ ˙ = [3 ¥ (3.14) 2 ¥ 6.25 ¥ 1028 ]1/ 3
Î V˚
= 1.23 ¥ 1010 per meter
Total kinetic energy of free electrons per unit volume at 0 K = (Average energy per electron at 0 K)
¥ (number of electrons per unit volume)
3 N 3
= EF ¥ = ¥ 5.75 ¥ 6.25 ¥ 1028 eV
5 V 5
= 21.56 1028 eV
E xamplE 7 Calculate the drift velocity of electrons in an aluminium wire of diameter 0.9 mm carrying
current of 6 A. Assume that 4.5 ¥ 1028 electrons/m3 are available for conduction.
d 0.9 ¥ 10 -3
Solution Given I = 6 A, n = 4.5 ¥ 1028 electrons/m3 and radius r = = = 4.5 ¥ 10 -4 m
2 2
I 6.0
Current density J= =
A p ¥ (4.5 ¥ 10-4) 2
6.0
=
3.14 ¥ (4.5 ¥10-4) 2
= 9.44 ¥ 106 A/ m 2
J 9.44 ¥ 106
and drift velocity vd = =
ne 4.5 ¥ 1028 ¥ 1.6 ¥ 10 -19
= 1.311 10 3 m / sec
E xamplE 8 The density of Cu is 8.92 ¥ 103 kg/m3 and its atomic weight is 63.5. Determine the current
density if the current of 5.0 A is maintained in Cu wire of radius 0.7 mm. Assuming that only one electron of
an atom takes part in conduction. Also calculate the drift velocity of electrons.
Solution Given
Atomic weight = 63.5 kg,
Density of copper = 8.92 ¥ 103 kg/ m3 , I = 5 A
Radius = 0.7 ¥ 10 -3 m
N 6.02 ¥ 1028 ¥ 8.92 ¥ 103
Radio =
V 63.5
= 8.456 ¥ 1030 electrons/m3
I I
Current density J = =
A pr2
5.0
= = 3.25 ¥ 106 m/sec
3.14 ¥ (0.7 ¥ 10-3) 2
J 3.25 ¥ 106
Drift velocity ( vd ) = =
ne 8.456 ¥ 1030 ¥ 1.6 ¥ 10-19
= 2.4 ¥ 10-6 m / sec
E xamplE 9 Using the data given, evaluate Fermi energy of the following alkali metals.
Li Na K
Density (g/cm3) r 0.534 0.971 0.86
Atomic Weight WA 6.939 22.99 39.202
h = 6.62 ¥ 10–34 J sec and m = 9.1 ¥ 10–31 kg
2 2/3 2 2/3
Solution Fermi energy EF =  ÈÍ3p 2 ◊ N ˘˙ = h2 ÈÍ3p 2 N ˘˙
2m Î V˚ 8p m Î V˚
N Nr
For Li =
V WA
N 6.023 ¥ 1026 ¥ 0.534 ¥ 103
= = 4.635 ¥ 1028 electrons/m3
V 6.939
(6.62 ¥ 10-34) 2
EF = ¥ [3 ¥ (3.14) 2 ¥ 4.635 ¥ 1028 ]2 / 3
8(3.14) 2 ¥ 9.1 ¥ 10-31
= 7.535 ¥ 10-19 J
= 4.71 eV
Similarly for Na,
N 6.023 ¥ 1026 ¥ 0.971 ¥ 103

= = 2.58 ¥ 1028 electrons/m3
V 22.99
(6.62 ¥ 10 -34) 2
EF = ¥ [3 ¥ (3.14) 2 ¥ 2.58 ¥ 1028 ]2 / 3
8 ¥ (3.14) 2 ¥ 9.1 ¥ 10 -31
EF = 5.032 ¥ 10 -19 J
or EF = 3.145 eV
and for K,
N 6.023 ¥ 1026 ¥ 0.86 ¥ 103
= = 1.321 ¥ 1028 electrons/m3
V 39.202
(6.62 ¥ 10-34) 2
EF = ¥ [3 ¥ (3.14) 2 ¥ 1.321 ¥ 1028]2 / 3
8 ¥ (3.14) 2 ¥ 9.1 ¥ 10-31
EF = 5.032 ¥ 10-19 J
EF = 3.145 eV
E xamplE 10 Calculate the energy difference between the ground state and first excited state for an electron
in one-dimensional box of length 10–10 m, V = 0 for 0 £ x £ a and V = • for 0 > x > a.
Solution Under the given conditions, the energy for the nth level is
2k 2 n2h2
En = = [where a is length of one-dimensional box]
2m 8ma 2
h2 È np ˘
so, E1 = ( n = 1) Íand k =
8ma 2 Î a ˙˚
4h 2
and E2 = ( n = 2)
8ma 2
The energy difference between first excited and ground state

3h 2 3 ¥ (6.62 ¥ 10-34) 2
DE = E2 - E1 = =
8ma 2
8 ¥ 9.1 ¥ 10-31 ¥ (10-10) 2
= 1.806 ¥ 10-17 J = 112.87 eV
Q.1 With the increase in temperature, the resistance of a metal

(a) remains constant (b) increases
(c) decreases (d) becomes zero
Q.2 Metals
(a) obey Ohm’s law
(b) have high electrical and thermal conditivities
(c) resistivity in proportional to temperature
(d) all of these
Q.3 Average kinetic energy ( E0 ) of a free electron gas at 0 K is
2 5 3
(a) EF (b) EF (c) EF (d) EF
5 3 5
Q.4 The density of states of electrons between the energy range E and E + dE is proportional to
(a) E1/2 (b) E2 (c) E (d) E3/2
Q.5 The phase space is a
(a) two dimensional space (b) one dimensional space
(c) three dimensional space (d) six dimensional space
Q.6 At low temperature, the resistivity of a metal is proportional to
(a) T 2 (b) T (c) T 5 (d) T1/2
Q.7 Which one of the following relations is correct for current density
1 ne
(a) J = nevd (b) J = (c) J = (d) J = neAvd
nevd vd
Q.8 In the presence of applied field, the average distance travelled by an electron between two successive
collisions is known as
(a) mean free path (b) drift velocity
(c) mobility of electron (d) none of these
Q.9 Which of the following is the correct form of Ohm’s law
(a) J = sE (b) J = s/E (c) J = sE2 (d) none of these
Q.10 Which one of the following relations is correct for the conductivity of metals
ne3 2m ne 2 t n
(a) s = 2mt (b) s = 2 (c) s = (d) s =
ne t 2m 2me 2 t
Q.11 The value of Fermi-distribution function at absolute zero (T = 0 K) is 1, i.e., F(E)=1, under the condition
(a) E > EF (b) E < EF (c) E = EF (d) E >> EF
Q.12 At any temperature T and for E = EF, the Fermi-distribution function becomes
1
(a) 0 (b) • (c) 1 (d) 2
Q.13 The free electron theory of metals was initiated by
(a) Pauli (b) Sommerfield
(c) Lorentz and Drude (d) Fermi-Dirac
Q.14 Which one of the following theory was developed by Lorentz and Drude
(a) quantum free electron theory (b) classical free electron theory
(c) zone theory (d) all of these
Q.15 Quantum theory of free electrons in metals explain:
(a) electrical conductivity and themionic emission
(b) specific heat & paramagnetism
(d) none of these
Q.1 What do you understand by free electron gas model of metals?

Q.2 What is the difference between electrical conductivity and thermal conductivity?
Q.3 Define and discuss Weidemann-Franz law.
Q.4 What are the limitations of free electron theory?
Q.5 Discuss the elements of classical free-electron theory. What were its limitations?
Q.6 What is the difference between classical and quantum theory of free electrons?
Q.7 Write expression for the Fermi energy at absolute zero.
Q.8 Write the expression for Fermi-Dirac distribution function.
Q.9 What are the drawbacks and successes of classical free electron theory?
Q.10 What is density of states? Discuss briefly.
Q.11 Write Richardson’s thermionic equation and explain every term occurring in this.
P ractice P roblems
General Questions
Q.1 What do you mean by free electron gas model of metals? Define free electron Fermi gas. Which
properties of solids are explained by free electron gas theory?
Q.2 Discuss the successes and failures of the free electron theory.
Q.3 Obtain an expression for the electrical conductivity of a metal on the basis of free electron theory.
Hence prove Ohm’s law.
Q.4 (a) Obtain an expression for thermal conductivity of a metal on the basis of free electron theory.
(b) State Weidemann-Franz law.
Q.5 Derive an expression for Fermi energy and density of states of a system.
Q.6 What is free electron theory of metals? Derive an expression for conductivity of metals on the basis of
Drude-Lorentz theory.
Q.7 Explain the quantum theory of free electrons in metals. Derive an expression for the Fermi-energy at
absolute zero.
Q.8 Discuss quantum theory of free-electrons and explain the following (a) Fermi level, (b) Density of
states (c) F-D distribution function.
Q.9 In terms of Fermi energy, calculate the kinetic energy at 0 K.
Q.10 Derive an expression for Fermi energy of free electrons. Discuss briefly the effect of temperature.
Q.11 Obtain an expression for energy levels in one dimensional free electron gas.
Q.12 What is Fermi gas? Does the Fermi energy of a metal depend upon the temperature?
Q.13 Considering the free electrons in a metal to form an electron gas obeying Fermi-Dirac statistics, obtain
Richardson’s equation for thermionic emission of electrons.
Q.14 State the difference between quantum and classical theories of free electron. Obtain Richardson-
Dushman equation of thermionic equation.
Q.15 Discuss the phenomenon of thermionic emission in metals. Obtain Richardson-Dushman equation for
the emission of current density.
Q.16 Derive the Richardson’s thermionic emission equation.
3
Q.17 Show that the kinetic energy of a three-dimensional gas of M free electrons at 0 K is E0 = EF .
2
(i) Fermi-Dirac distribution function.
(ii) Density of states in one-dimension.
(iii) Energy levels and wave function of free electrons in a box.
Band Theory of Solids and

Photoconductivity
18
Learning Objectives
LO 1 Understand Kronig-Penney model, LO 4 Explain electron configuration in
energy E versus k diagram and one- and conduction band, hole concentration
two-dimensional Brillouin zones in valence band, Fermi level of intrinsic
LO 2 Explain effective mass of an electron and extrinsic semiconductor
and deviation of electron behaviour in LO 5 Illustrate Hall effect, Hall voltage and
crystal lattice Hall coefficient
LO 3 Differentiate between Insulators, LO 6 Learn simple model of photoconductor
n-type, p-type semiconductors and and gain factor
conductors (metals) LO 7 Discuss effects of traps, applications of
photoconductivity
Introduction
A solid contains an enormous number of atoms packed closely together. When N atoms of the solid are
well separated, then these atoms lead to N-fold degenerate levels of the solid. As the atoms approach
one another to form a solid, i.e., their separation reduces, a continuously increasing interaction occurs
between them. This causes each of the levels to split into N distinct levels. It is the separation distance
(say r) which specifies the amount of overlap that causes the splitting. Since a solid contains about 1023
atoms per mole, i.e., N is very large, the splitted energy levels become so numerous and close together
that they form an almost continuous energy band.
The amount of splitting is different for different energy levels. For example, the lower energy levels are
found to spread or split less than the higher levels. It means the lowest levels remain almost unsplit. The
reason is that the electrons in lower levels are the ones which are in inner subshells of the atoms. So they
are not significantly influenced by the presence of nearby atoms. Since the potential barriers between the
atoms are for them relatively high and wide, these electrons are localised in particular atoms, even when
r is small. However, the electrons in the higher levels are the valence electrons and are not localised at all
Band Theory of Solids and Photoconductivity 655
for small r but they become part of the whole system. From the quantum point of view, the wave functions
of the valence electrons overlap and the overlapping of their wave functions results in splitting or spreading
of their energy levels.
The band formation of the higher energy levels of
sodium, whose ground state atomic configuration
is 1s2 2s2 2p6 3s1, is shown in Fig. 18.1. In the figure,
the dashed and vertical line indicates the observed
interatomic separation in the solid sodium. It is clear
from the figure that the bands overlap when the atomic
separation decreases. This figure also shows that the
allowed band corresponding to inner subshells, for
example 2p in sodium, are extremely narrow and
does not begin to split until the interatomic distance
r becomes less than the value actually found in the
crystal. As we move towards the higher energy states,
the energy of the electrons become larger and also the
region in which they can move becomes wider. Since
they are also affected more by the nearby ions, it is
seen that the bands become progressively wider for the Figure 18.1
outer occupied subshells and also for the unoccupied
subshells of the atoms in its ground state. Therefore, with the increase of energy the successive allowed
bands become wider and overlap each other in energy.
It is clear from the above discussion that the energy bands in a solid correspond to energy levels in an
atom. Therefore, an electron in a solid can occupy only energy that falls within these energy bands.
The overlapping of the bands depends on the structure of the solid. If the bands do not overlap, then
the intervals between them represents energies which the electrons in the solid cannot occupy. These
intervals are called forbidden bands or energy gaps. However, if the adjacent bands in the solid overlap,
then the electrons possess a continuous distribution of allowed energies.
18.1 KrONiG-PENNEy mODEl LO1

Free electron model ignores the effects those arise when the electrons interact with crystal lattice. However,
now we consider this by making some general remarks about the effect of the periodic variation in the potential.
Due to this periodicity in potential for an infinitely long lattice, the wave function does not remain sinusoidal
travelling waves of constant amplitude but now they include the lattice periodicity in their amplitudes, and
electrons may be scattered by the lattice. When the deBroglie wavelength of the electron corresponds to a
periodicity in the spacing of the ions, the electron interacts strongly with the lattice. This situation is the same
as an electromagnetic wave suffers Bragg’s reflection, when the Bragg’s condition is satisfied.
In order to find the allowed energies of electrons in solids, we consider the effect of formation of a solid
when the individual constituent atoms are brought together. We solve the Schröedinger equation for periodic
potential seen by an electron in a crystal lattice. We also consider that the periodic potential is a succession of
rectangular wells and barriers. The solution of Schröedinger equation is a sinusoidal wave in certain energy
ranges, i.e., allowed states, and real decaying exponential wave in the other ranges, i.e., the forbidden bands.
For this purpose, here we present only qualitative approach.
It is found that the potential is not constant but varies periodically. The effect of periodicity is to change
the free particle travelling wave eigen function. Therefore, the travelling wave eigen function has a varying
amplitude which changes with the period of the lattice. If we consider that the space periodicity is a (Fig.
18.2), then according to Bloch, the eigen function for one-dimensional system has the form
y(x) = uk(x)eikx
V(x)
V0
–b 0 a
(a + b)
Figure 18.2
As is clear, this is different from the free travelling wave function y(x) = Aeikx ◊ uk(x) is the periodic function
with the periodicity a of the periodic potential, i.e.,
uk(x) = uk(x + a)
In general,
uk(x) = uk(x + na)
where n is an integer. Hence, with the effect of periodicity, the complete wave function is
y(x, t) = uk(x)ei(kx – wt) (i)
2p
In the above equation, the exponential term indicates a wave of wavelength l = which travels along +x
k
direction if k is positive and it moves along –x direction if the value of k is negative.
The exact form of the function uk (x) depends on the particular potential assumed and the value of k.
In 1930, Kronig and Penney proposed a one-dimensional model for the shape of rectangular potential wells
and barriers having the lattice periodicity, as shown in Fig. 18.2. Each well represents an approximation to
the potential produced by one ion. In the region such as 0 < x < a, the potential energy is assumed to be zero
while in the region –b < x < 0 or a < x < (a + b), the potential energy is taken as V0. The relevant Schröedinger
equations for these two regions are
d 2y È 2m ˘
+ Í 2 ˙ Ey = 0 (ii) [0 < x < a ]
2
dx Î ˚
d 2y È 2m ˘
+ Í 2 ˙ ( E - V0 )y = 0 (iii) [–b < x < 0 ]
2
dx Î ˚
The electron of not too high energy is practically Single Potential Well Periodic Array of Wells
bound within one of the wells that are deep and widely
Energy
spaced. So the lower energy eigen values are those of
a single well. However, for the wells those are closer
together the eigen function can penetrate the potential 2V0
barriers more easily. Because of this, spreading of
previously single energy level into a band of energy
levels takes place. The band becomes wider with the
decrease in the separation of the wells. Under the limit
of zero barrier thickness, we obtain an infinitely wide
single well in which all energies are allowed. So the
present case is reduced to the free electron model. The V0
comparison between the allowed energies of a single
well and an array of wells (Kronig-Penney model) is 0.9
shown in the Fig. 18.3. In this figure, we have assumed
b = a/16 and the well strength as 2mV0a2/2 = 121. It 0.5
is clear from the figure that each band corresponds to
a single energy level of the single well. The forbidden 0.2
bands appear even for energies E > V0. 0.06
0
Here we will solve the Schröedinger wave equation
for electron for Kronig-Penney potential under the Figure 18.3
dy
condition that y and are continuous at the boundaries of the well. A complicated expression for the
dx
allowed energies in terms of k shows that gaps in energy are obtained at values such that
p 2p 3p
k=± ,± ,± ,º (iv)
a a a
The solution of the Schröedinger wave equation for free-electrons results in the energy values given by
h2 k 2 2 k 2
E= 2
= (v)
8p m 2m
18.1.1 Energy E versus Wave Number k Diagram

From Eq. (v), it is clear that the relation between E and k is parabolic. The parabolic relation between E and k,
valid in case of free electrons, is therefore, interrupted at different values of k, as shown in Fig. 18.4. It means
the energies corresponding to the values of k given by Eq. (iv) are not permitted for electrons in the crystal.
Thus, the energies of electrons are divided into forbidden and allowed bands (Fig. 18.4).
Figure 18.4
The occurrence of the gaps can be understood on the basis of Bragg’s condition for the diffraction, given as
2a sin q = nl n =1, 2, 3, ….. (vi)
where a is spacing between the ions of the lattice and q is the angle of incidence.
Eq. (vi) can be written as
2a = nl (for q = 90º)
2p
or 2a = n
k
np
or k= (vii)
a
p 2p 3p
or k =± ,± ,± ,º
a a a
We have put ± signs because the incident wave can travel along +x-axis as well as along –x-axis. At all these
values of k the gaps in energy occur, as shown in Fig. 18.4.
The waves corresponding to values of k satisfying the Bragg’s condition are reflected and resulted in standing
waves. On each subsequent Bragg reflection, the direction in which the wave is travelling is reversed again.
p
The eigen function of incident and corresponding reflected waves for k = ± n are therefore ei(p/a)x and
a
e–i(p/a)x. These two eigen functions can be combined in two different ways to give total eigen function
y1 = ei(p/a)x + e–i(p/a)x = cos(p/a)x
y2 = ei(p/a)x + e–i(p/a)x = sin(p/a)x

or y1 µ cos(p/a)x (viii)
and y2 µ sin(p/a)x (ix)
Hence, the two standing waves are obtained. The probability density curves for these two stationary waves,
i.e., |y1|2 and |y2|2, are shown in Fig. 18.5. From this figure and Eq. (viii) it is clear that the value of |y1|2 is
maximum at the positions of positive ions (i.e., x = 0, ±a, ±2a,...). The value of |y2|2 is maximum in between
the position of positive ions. From Fig. 18.2, it is evident that the potential energy of an electron is maximum
between the ions and minimum at the positions of the ions. So an electron can have two different values of
p
energies, i.e., E1 and E2 for k = corresponding to the two standing waves y1 and y2. Hence, no electron
a
can have any energy between E1 and E2. This phenomenon creates a difference in energy (E1 ~ E2) which is
known as energy gap.
�Y�2
�Y1�2 �Y2�2
–3a –2a –a 0 a 2a 3a
Figure 18.5
18.2 ONE-aND TWO-DimENsiONal BrillOuiN ZONEs LO1

In Kronig-Penney model, we have seen that the discontinuities in energy occurs when the wave number
k satisfies the condition k = np/a, where n takes the values ±1, ±2, ±3,... etc. The graph between the total
energy E and wave number k is shown in Fig. 18.4. It is clear from the figure that an electron has allowed
energy values in the region between –p/a to +p/a. This region is called the first Brillouin zone. As discussed
earlier, there is a discontinuity of gap in the energy values after this allowed energy value. This gap is called
forbidden gap or forbidden zone. Again there is another allowed energy zone, which is observed after this
forbidden gap and is extended from –p /a to –2p /a and p/a to 2p /a. This zone is called second Brillouin zone.
Similarly, the other higher order Brillouin zones can be defined.
The first two Brillouin zones in one-dimensional case are shown in Fig. 18.6. We can extend the concept
of the Brillouin zones to two-dimension by considering that the electron is moving in the two-dimensional
square lattice. In this case, the wave number k has the two components, i.e., along the x-axis and y-axis.
Let us represent them as kx and ky, respectively. The two-dimensional Brillouin zones are shown in Fig.
18.7. It can be seen that kx = ky = ±p/a limits the first Brillouin zone. In the figure, the first Brillouin zone
is represented by a square passing through the points A, B, C and D. The second Brillouin zone for a two-
dimensional lattice is represented by a square passing through E, F, G and H.
Second Brillouin Zone
–3p/a –2p/a –p/a p/a 2p/a 3p/a
K
First Brillouin Zone
Figure 18.6
It is clear that the boundary for the second zone is

given by kx=±2p/a and ky=±2p/a. Thus the region
between the squares ABCD and EFGH is called the
second Brillouin zone.
18.3 EffEcTivE mass Of LO2

aN ElEcTrON
The electrons in a crystal are not completely free
but interact with the crystal lattice. As a result, their
behaviour towards external forces is different from
that of a free electron. The deviation of electron
behaviour in the crystal lattice from the free electron
behaviour can be taken into account simply by
considering the electron to have an altered value of
mass called the ‘effective mass’ m* rather than its
mass m* in free space, which is different from the
mass m of the electron in free space. The effective Figure 18.7
mass m* depends on the nature of crystal lattice and varies with the direction of motion of the free electron
in the lattice.
Suppose an electron is moving along the x-axis in a crystal in the presence of an external electric field E¢.
So it experiences a force eE¢. If the electron gains velocity v over a distance dx in time dt under the action of
this force, then
Work done = dE=eE¢ dx= eE¢ v dt È dx ˘
ÍÎ v = dt ˙˚
As we know that the velocity v of a particle (electron) is the same as the ‘group velocity’ ÊÁ vg =
dw ˆ
Ë ˜ of the
dk ¯
de Broglie waves associated with the particle (v = vg). Thus, we can write the work done in terms of vg as
dE = eE¢ vg dt (i)
According to Einstein’s de Broglie relation
h
E = hn = w [ w = 2pn]
2p
By differentiating it, we get
h h dw (ii)
dE = dw = dk
2p 2p dk
dw
= vg
 dk
h
 dE = vg dk (iii)
2p
By comparing Eq. (iii) with Eq. (i), we get

h
eE ¢dt = vg dk
2p
dk 2p
or = eE ¢ (iv)
dt h
Let us write the group velocity vg in terms of energy E from Eq. (ii)
d w 2p dE
vg = = (v)
dk h dk
By differentiating the above equation w.r.t. time, we get
dvg 2p d 2 E 2p d 2 E dk
= =
dt h dtdk h (dk ) 2 dt
dvg 2p d 2 E 2p
or = eE ¢ (By using Eq. (iv)) (vi)
dt h (dk ) 2 h
Employing vg = v again, this can be written as
dv Ê 4p 2 d 2 E ˆ (vii)
=Á ˜ eE ¢
dt Ë h 2 dk 2 ¯
dv
This equation connects the force eE¢ on the electron with the acceleration through the proportionality
dt
Ê 4p 2 d 2 E ˆ
factor Á 2 .
Ë h dk 2 ˜¯
Since F = ma, a = F/m (viii)
A comparison of Eq. (viii) with Eq. (vii) yields
1 4p 2 d 2 E
=
m* h 2 dk 2
1
The quantity is the reciprocal of the effective mass of the electron in the crystal lattice.
m*
18.4 DisTiNcTiON BETWEEN iNsulaTOrs, sEmicONDucTOrs aND LO3

cONDucTOrs (mETals)
The formation of bands in solids has already been discussed and it was shown that there is an energy gap,
called forbidden band, representing energies which the electrons cannot occupy. Based on this energy gap
and the conduction, the solids are classified into different categories named as insulators, semiconductors
and conductors.
18.4.1 insulators
For these types of solids, the band formation is like the one shown in Fig. 18.8a. In this case, the forbidden gap
between the highest filled band (valence band) and the lowest empty band (conduction band) is very wide; it
is about 3 eV to 6 eV. It is seen that a very few electrons from the filled band reach the empty band, even if
we thermally excite them or apply an electric field to them. Moreover, Pauli exclusion principle restricts the
electrons for moving about in the filled band. For this reason, a free electron current cannot be obtained and
the solids of this type are poor conductors of electricity. This class of solid is known as insulators. Diamond,
quartz, and most covalent and ionic solids like ZnO and AgCl are the examples of insulators.
Empty 3s
Empty 3s
Forbidden
About 3 eV – 6 eV
Gap
About 0.1 eV – 1 eV Forbidden Gap
Filled 2p Filled 2p
Filled 2s Filled 2s
1s 1s
Figure 18.8a Figure 18.8b
18.4.2 semiconductors
For these types of solids the band formation is like the one shown in Fig. 18.8b. In this case, the forbidden gap
between the highest filled band (valence band) and the lowest empty band (conduction band) is very narrow;
it is about 0.1 eV to 1 eV. Under this situation we can easily move the electrons from the highest filled band
to the empty band. This can be achieved by thermal excitation or also by applying an electric field. For this
reason, a free electron current can be obtained as a few electrons are available in the empty band. This class
of solids is known as semiconductors. Silicon and germanium are the examples of semiconductors.
In semiconductors, there also exists another mechanism that causes the generation of electric current.
Actually there are vacancies or the empty places left behind when the electron moves, which remain near the
top of the uppermost filled band. These vacancies are called holes. The holes behave as positive electrons
and can contribute to the generation of electric current. This is possible as the electron below the hole may
gain enough energy to jump and occupy the hole due to the applied electric field. With such successive jumps
of the electrons, the hole moves towards the lower energy state and contribute to the generation of electric
current.
The semiconductors are mainly of two types, defined below.

18.4.2.1 n-type Semiconductors
The n-type semiconductors are the ones in which the electron conduction (negative) exceeds the hole
conduction (positive). In such semiconductors, the donor impurity predominates. This can be understood if
we introduce a small amount of phosphorous (P) or arsenic (As), i.e., an element of fifth group of the periodic
table, as an impurity into a crystal of silicon (Si) or germanium (Ge). This addition of P or As means replacing
an atom of Si or Ge at a lattice site by an atom of the impurity. Atoms of fifth group elements have five valence
electrons whereas Si or Ge has four valence electrons. So four electrons of P or As form covalent bonds with
the electrons of the atoms of Si or Ge. However, the fifth electron remains only very weakly bound to the
P or As atom by electrostatic forces and this cannot be accommodated in the already filled original valence
band. So, it occupies a discrete energy level which is just below the conduction band (with only a few tenths
of an eV). Hence these extra electrons jump easily into the conduction band and contribute to the electric
conductivity in addition to the electron hole pairs produced by thermal excitation of the pure semiconductor.
This way the number of electrons sits more than holes to serve as charge carriers.
18.4.2.2 p-type Semiconductors
The p-type semiconductors are the one in which the hole conduction (positive) exceeds the electron
conduction (negative). In such semiconductors, the acceptor impurity predominates. This can be understood
if we introduce a small amount of Al, Ga or In, i.e., an element of third group of the periodic table, as an
impurity into a crystal of silicon (Si) or germanium (Ge). Atoms of third group elements have three valence
electrons whereas Si or Ge has four valence electrons. So three electrons of Al or Ga form covalent bonds with
the electrons of the atoms of Si or Ge. However, the fourth available electron of the semiconductor lacks an
electron with which it can form a bond. This is equivalent to as if a vacancy or hole has been created at the site
of the impurity atom. Hence, the impurity atoms introduce vacant discrete energy levels very near the top of
completely filled valence band of Si or
Ge. So these extra holes move from an
impurity atom. These holes behave as
positive charge carriers and are available
in excess. Since the crystals of this type
have an excess of positive charge carriers,
they are called positive semiconductors
or p-type semiconductors.
18.4.3 conductors or metals
For these types of solids, the band
formation is like the one shown in
Fig. 18.8c. In this case, the valence
band is either partially filled or the next
allowed empty band overlaps with the
filled band. In both the cases, there are
unoccupied states for electrons in the
uppermost band. So these electrons are
available to generate the current. This
class of solids is known as conductors. Figure 18.8c
The conductors offer a low resistance to
the passage of an electric current. Silver, copper, iron, aluminium, etc. are the examples of conductors or
metals.
18.5 iNTriNsic sEmicONDucTOr LO4

The semiconductors in which the transformation of electrons to the conduction band and the generation of holes
in the valence band are achieved purely by thermal excitation are called intrinsic semiconductors. It means this
effect is temperature dependent and produces equal numbers of electron and hole carriers. The electrons and
holes are called intrinsic charge carriers and the resulting conductivity is known as intrinsic conductivity.
18.5.1 Electron concentration in conduction Band

For this we assume the following.
(i) Electrons in the conduction band behave as free particles with an effective mass m.
(ii) The number of conduction electrons per cubic meter whose energies lie between E and E + dE is
given by
dnC = N(E) f (E) dE (i)
where N(E) is the density of states at the bottom of the conduction band and is
given by as per quantum mechanics.
4p
N (E) = (2m)3/2 ( E - EC )1/2 (ii)
h3
Here EC is the energy at the bottom of the conduction band.
(iii) The Fermi-Dirac probability function which specifies the fraction of all states at energy E occupied
under condition of thermal equilibrium is given by
1
f (E) = (iii)
1 + e( E - EF ) / kT
Here k is the Boltzmann constant, T is the temperature in Kelvin and EF is the Fermi level or
characteristic energy for the crystal.
Now since the electrons in the conduction band may have energies lying between EC and •, the
concentration of electrons in the conduction band is given by
•
nC = Ú N ( E ) f ( E )dE (iv)
EC
where EC is the energy at the bottom of the conduction band.

Combining Eqs. (ii), (iii) and (iv), we get the density of electrons in the conduction band as
( E - EC )1/2
•
4p (v)
nC =
h3
(2m)3/2 Ú 1 + e( E - E F ) / kT
dE
EC
For E ≥ EC and E – EF >> kT, l in the denomination can be neglected and this equation reduces to
•
4p 1/2 ( EF - E ) / kT
nC =
h3
(2m)3/2 Ú ( E - EC ) e dE
EC
•
4p
nC = 3
(2m)3/2 e( EF - EC ) / kT Ú ( E - EC )
1/2 ( EC - E ) / kT
e dE
h EC
Substituting (E – EC) / kT = x, so that dE = kT dx, we get

•
4p
nC = (2m)3/2 e( EF - EC ) / kT Ú x1/2 (kT )1/2 e - x kT dx
h3 0
•
4p
= (2m)3/2 (kT )3/2 e( EF - EC ) / kT Ú x1/2 e - x dx
h3 0
4p p È •
p ˘
= 3
(2m)3/2 (kT )3/2 e( EF - EC ) / kT Í\ Úx
1/2 - x
e dx = ˙
h 2 ÍÎ 0
2 ˙˚
3/2
È 2p mkT ˘
nC = 2 Í e( EF - EC ) / kT (vi)
Î h 2 ˙˚
This relation gives the density or concentration of electrons in the conduction band of an intrinsic
semiconductor. Please note that here m is the effective mass of the electron.
18.5.2 Hole concentration in valence Band

By a similar method as adopted above, we can get an expression for the concentration of holes in the valence
band. Since a hole signifies a vacancy created by removal of an electron, i.e., an empty energy level, the Fermi
function for a hole is 1 – f (E). Here f (E) represents the probability that the level is occupied by an electron.
Along with the expression of f (E), the factor 1 – f (E) becomes
1
1 - f (E) = 1 - ( E - EF ) / kT
1+ e
= 1 - [1 + e( E - EF ) / kT ]-1
= 1 - [1 - e( E - EF ) / kT ] for ( E - EF )/ kT << 1
( E - EF ) / kT
=1-1+ e
1 - f ( E ) = e( E - EF ) / kT
For the top of the valence band (the maximum energy), the density of states is given by
4p
N (E) = 3
(2mh )3/2 ( Ev - E )1/2
h
Here mh is the effective mass of holes near the top of the valence band, where the energy is EV. With the above
relation, the density of holes in the valence band is calculated as
EV
nh = Ú N ( E )[1 - f ( E )]dE
-•
Ev
4p 1/2 ( E - EF )/ kT
=
h3
(2mh )3/2 Ú ( Ev - E ) e dE
-•
EV
4p
= (2mh )3/2 e( EV - EF )/ kT Ú ( EV - E )
1/2 ( E - EV )/ kT
e dx
h3 -•
0
4p
= 3
(2mh )3/2 e( EV - EF )/ kT Ú x1/2 (kT )1/2 e - x kTdx
h +•
E -E
where we have substituted V so that dE = – kT dx. Now
kT
•
4p
nh = (2mh ) 2/3 (kT )3/2 e( EV - EF )/ kT Ú x1/2 e - x dx
h3 0
3/2
È 2p mh kT ˘
nh = 2 Í ˙ e( EV - EF )/ kT (vii)
Î h2 ˚
18.5.3 intrinsic concentration of charge carriers
Combining Eqs. (vi) and (vii), we get the following expression for the product of electron-hole concentration
3
È 2p kT ˘
nC nh = np = 4 Í 2 ˙ (mmh )3/2 e( EV - EC )/ kT
Î h ˚
- E g / kT
nC nh = AT 3e (viii)
where Eg = EC – EV is the width of forbidden energy gap between conduction and valence bands and
32p 3 k 3
A= (mmh )3/2 is a constant. In most of the cases, nC is written as n and nhas p only.
h6
Eq. (viii) shows that the product of holes and electron densities depends on the temperature T and the
forbidden energy gap Eg, but is independent of the Fermi level EF. Thus the product of electron and hole
concentrations, for a given material, is constant at a given temperature. If an impurity is added to increase
n, there will be a corresponding decrease in p such that the product np remains a constant. Since for an
intrinsic semiconductor, n = p = ni, we arrive at an important relationship, called the law of action
- E g / kT
np = ni2 = AT 3e (ix)
where ni is called the intrinsic density of either carrier. Equation (ix) is true for a semiconductor regardless
of donor or acceptor concentrations.
18.5.4 Energy Band Diagram and fermi level

In an intrinsic semiconductor, the electrons and holes are always generated in pairs, i.e., n = p = ni. Substituting
the values of n and p from Eqs. (vi) and (vii) of the previous sections, we get
3/2 3/2
È 2p mkT ˘ È 2p mh kT ˘
2Í e( EF - EC ) / kT = 2 Í ˙ e( EV - EF ) / kT
Î h 2 ˙˚ Î h2 ˚
3/2
Êm ˆ
or e(2 EF - EC - EV ) / kT = Á h ˜
Ë m¯
Taking the logarithm of both sides

2 EF - EC - EV 3 Ê mh ˆ
= ln Á ˜
kT 2 Ë m¯
3 Êm ˆ
or 2 EF - EC - EV = kT ln Á h ˜
2 Ë m¯
E + EV 3 Êm ˆ
or EF = C + kT ln Á h ˜
2 4 Ë m¯
If the effective mass of hole and a free electron is the same i.e., mh = m, then
EC + EV
EF =
2
This shows that the Fermi level EF lies exactly in the centre of
the forbidden energy gap Eg as depicted in Fig. 18.9. The Fermi
level can also be defined as the energy level at which there is a 0.5
probability of finding an electron. It depends on the distribution of
energy levels and the number of electrons available. Figure 18.9
18.6 ExTriNsic sEmicONDucTOr LO4

The conductivity of an intrinsic semiconductor can be increased significantly by adding certain impurities to
it. By doing so, we get impurity semiconductor which is also known as extrinsic semiconductor. In extrinsic
semiconductors, current carriers are introduced by donor and acceptor impurities with locked energy levels
near the top or bottom of the forbidden gap.
18.6.1 Energy Band Diagram and fermi level

(i) n-type Extrinsic Semiconductor
When a small amount of pentavalent impurity is added to the crystal, it creates extra electrons without adding
any new holes. This impurity introduces new energy levels into the energy band picture. The location of these
new levels is slightly below the bottom of the conduction band for intrinsic semiconductor. The width of the
gap for germanium was previously stated as 0.72 eV. The energy required to move an electron from a donor
impurity into the conduction band is of the order of 0.01 eV, and since at normal ambient temperature the
thermal energy is considered to be about 0.02 eV, it is concluded that almost all the electrons are detached
from the donor atoms and have conduction band energies. In the case of silicon doped with donor impurities,
the energy required to move an electron from donor impurity into the conduction band is of the order of
0.05 eV. The energy band diagram for an n-type semiconductor is shown in Fig. 18.10. Here Eg represents the
energy level corresponding to donor impurities.
In intrinsic semiconductor, Fermi level lies in the middle of the

forbidden energy Eg indicating equal concentrations of free electrons
and holes. When a donor type impurity is added to the crystal, then
if we assume that all the donor atoms are ionised, the donor electrons
will occupy the states near the bottom of the conduction band. Hence,
it will be more difficult for the electrons from the valence band to cross
the energy gap by thermal agitation. Consequently, the number of
holes of the valence band is decreased. Since Fermi level is a measure Figure 18.10
of the probability of occupancy of the allowed energy states, EF for
n-type semiconductors must move closer to the conduction band, as shown in Fig. 18.10. At usual temperatures
all the donor levels will be fully activated and the donor atoms will be ionised. It means the density of electrons in
the conduction band will be approximately equal to the density of donor atoms, i.e., n ª Nd (Nd being the density
of donor atoms). Then from Eq. (vi), we have
3/2
È 2p mkT ˘
n = Nd = 2 Í e( EF - EC ) / kT
Î h 2 ˙˚
= N c e( EF - EC ) / kT
3/2
È 2p mkT ˘
where N c = 2 Í = constant
Î h 2 ˙˚
Nc
\ = e - ( EF - EC ) / kT
Nd
Taking the logarithm of both the sides
Nc E - EC
ln =- F
Nd kT
Nc
or EF = EC - kT ln
Nd
It shows that the Fermi level lies below the bottom of the conduction band, as shown in Fig. 18.10.
(ii) p-type Extrinsic Semiconductor

When an acceptor-type impurity is added, it also modifies
the energy level diagram of semiconductor and makes the
conduction easier. The presence of impurity creates new
energy levels which are in the gap in the neighbourhood of
the top of valence band of energies, as shown in Fig. 18.11.
Ambient temperature results in ionisation of most acceptor
atoms and thus an apparent movement of holes takes place
from the acceptor levels to the valence band. The energies
for holes are highest near the valence bond and decrease Figure 18.11
vertically upward in the energy level diagram. Alternatively, one may say that electrons are accepted by the
acceptors and these electrons are supplied from the valence band, thus leaving a preponderance of holes in
the valence band.
The energy band diagram for a p-type semiconductor is shown in Fig. 18.11 where EA represents the energy
level corresponding to the acceptor impurity. When an intrinsic semiconductor is doped with acceptor type
impurity, the concentration of holes in the valence band is more than the concentration of electrons in the
conduction band and the Fermi level shifts towards the valence band, as shown in Fig. 18.11. The acceptor
level lies immediately above the Fermi level.
If we assume that there are only acceptor atoms present and that these are all ionised, we have p = Na. Then
from Eq. (vii), we get
3/2
È 2p mh kT ˘
p = Na = 2 Í ˙ e( EV - EF ) / kT
Î h2 ˚
= N v e( EV - EF ) / kT
3/2
È 2p mh kT ˘
where N v = 2 Í ˙ = constant
Î h2 ˚
Nv
\ = e - ( EV - EF ) / kT
Na
Taking logarithm of both the sides,
Nv E - EF
ln =- V
Na kT
NV
or EF = EV + kT ln
Na
It shows that the Fermi level lies above the top of valence band, as shown in Fig. 18.11.
18.6.2 Effect of Temperature
Let us see what happens if we increase the temperature of an n-type semiconductor. Since all the donors have
already donated their free electrons at room temperature, the additional thermal energy will only increase the
generation of electron-hole pairs. Thus the concentration of minority charge carriers increases. A temperature
is ultimately reached when the number of covalent bonds broken is very large such that the number of holes
and electrons is almost equal. The extrinsic semiconductor then behaves like an intrinsic semiconductor,
although its conductivity is higher. This critical temperature is 85oC for germanium and 20oC for silicon. The
same arrangement can be put forward for the p-type semiconductor. Thus with an increase in the temperature
of an extrinsic (impurity) semiconductor, it behaves almost intrinsically.
18.7 Hall EffEcT LO5

If a current carrying conductor is placed in a transverse magnetic field, a potential is developed in the
conductor in the direction perpendicular to both the current and magnetic field. This phenomenon is known
as Hall effect. It was discovered by Hall in 1879.
Let us consider a rectangular strip carrying current along x-axis (with electron flow towards x-axis) and
magnetic field of strength B is applied along z-axis, as shown in Fig. 18.12. The force on electrons would be
exerted due to the effect of magnetic field in negative y-direction. Therefore, the face MN collects a negative
charge and the face PQ collects positive ions. This separation of charges sets up an electrostatic field inside
B
Z
Y
Jx
Q P
X
M N
(a)
Ea
Q P
d Jx
M N
(b)
Q P
Ey
M N
(c)
Figure 18.12
the conductor in y-direction. The accumulation of charge on the surfaces of the specimen continues until the
force on moving charges due to electric field associated with the accumulated charge itself is large enough
to cancel the force exerted by the magnetic field. Ultimately a steady state is reached in which the net force
on the moving charges in y-direction vanishes and the electron can again move freely along the conductor. In
stationary state the value of E is denoted by EH and is called Hall electric field.
18.7.1 Hall voltage and Hall coefficient

We can calculate Hall voltage and Hall coefficient by using the following approach.

The force on the charge carriers = qE H
 
The force on charge carriers due to magnetic field B = q(v d ¥ B)
At the stationary state both these forces balance each other. Hence
     
qE H = q(v d ¥ B) or E H = v d ¥ B (i)
Thus by measuring EH and B, the drift velocity vd can be determined both in magnitude and direction. If we
talk about magnitude only, then Eq. (i) gives
EH = vd B [\ vd and B are at right angle] (ii)
The current density J in terms of the drift velocity vd is written as
J = – nevd [n = number of charge carriers per unit volume] (iii)
From Eqs. (ii) and (iii), we get
JB
EH = - (iv)
ne
The above relation shows that
EH µ B
If d be the width of the strip, then the transverse Hall electric field EH can be related to the Hall potential
difference VH as
VH
EH =
d
JBd
or VH = EH d = - (v)
ne
VH is also known as Hall Voltage.
1
The coefficient of proportionality - is called Hall coefficient and is denoted by RH. It is given below
ne
1
RH = – (vi)
ne
Hall coefficient is negative if the charge carriers are electrons and it will be positive if the charge carriers are holes.
From the above equation, we can develop a relation between the Hall electric field EH and the electric field Ex,
1
which causes the current to flow in the conductor. We can write the current density J = s Ex = Ex where s
is the conductivity and r is the resistivity. Then from Eq. (iv) r
JB - Ex B /r
EH = - =
ne ne
Ex B
=-
enr
EH -B
=
E x ner
18.8 PHOTOcONDucTiviTy LO6

Photoconductivity is an electrical phenomenon in which a material becomes more conductive due to the ab-
sorption of light radiation. When the energy of the incident radiation is higher than the energy gap Eg between
conduction band and valence band, the electrons and holes raise the electrical conductivity. All the photo-
conductors are not the intrinsic type, but imperfections also play an important role. If impurities are present,
electron and hole pairs are produced even if the energy of incident photon is below the threshold. The role of
impurities is important to understand the experimental facts of photoconductivity. The presence of impurities
produces discrete energy levels in the forbidden gap and they are often known as traps.
18.9 simPlE mODEl Of PHOTOcONDucTOr LO6
We shall discuss a simple model of photoconductor in order to understand the photoconductivity. When light
radiations fall on the crystal specimen, electron-hole pairs are produced throughout the volume of the crystal,
as shown in Fig. 18.13.
The recombination process occurs by direct annihilation of electrons Light Radiation

with holes. We further suppose that as many electrons leave the crystal
at one electrode the same number of electrons enters into the crystal
from the opposite electrode. This model is highly hypothetical one. The
production rate of electron-hole pairs is –
+ – + +
dn
= L - Ane nh = L - An 2 (i) [ne = nh = n] –
+ –
–
dt + +
–
where L is the number of photons absorbed by the specimen, ne and nh
are the number of electrons and holes, respectively, and A is known as Figure 18.13
recombination coefficient. The term Anenh shows the recombination
rate and it will be proportional to the rate of production of electron-hole pairs.
In the steady state,
dn
=0
dt
so that, Eq. (i) becomes
L – An02 = 0
1/ 2
Ê Lˆ
or n0 = Á ˜ (ii)
Ë A¯
Here n0 is the number of electrons in the steady state.

The corresponding conductivity is given by
1/ 2
ÈL˘ (iii)
s = n0 eme = Í ˙ eme
Î A˚
where, me is the mobility of the electron.
Hence, the photocurrent density will be
J = sE
ÈL˘
1/ 2
ÈV ˘ È V˘
J =Í ˙ eme Í ˙ ÍÎ\ E = d ˙˚
Î A˚ Îd ˚
n0 emeV
or J= (iv)
d
where V is the voltage across the specimen, and d is the thickness of the specimen. If the light is switched off
suddenly, then L = 0, by which Eq. (i), becomes
dn dn
= - An 2 or = - Adt
dt n2
On integration, we get
dn n -2 + 1
Ú n2 = - AÚ dt or
-2 + 1
= - At - constant
1 1
or = At + (constant ) (v)
n n0
since at t = 0, n = n0 (numbers of electrons at time t = 0). If light is turned off, i.e., the light falling on the
1
specimen is stopped and generation of the electrons stops so that the electron concentration drops to n0 in
time t , then 2
0
1 1
= At0 +
n0 /2 n0
1 1 1
or t0 = = 1/ 2
= = ( AL ) -1/ 2
An0 A( L / A) ( AL )1/ 2
1/ 2
1 L Ê Lˆ 1 n0 1 n
or t0 = ¥ =Á ˜ = or t0 = = ( AL ) -1/2 = 0 (vi)
( AL ) 1/ 2
L Ë A¯ L L An0 L
s
since s = n0 eme or n0 = Eq. (vi) becomes
em e
n0 s
t0 = = (vii)
L eme L
where time t0 is known as the response time. From equation (vii) it is clear that the response time is directly
proportional to the photoconductivity at a given light level L.
18.9.1 sensitivity or Gain factor

Sensitivity is the ratio of number of carriers crossing the specimen to the number of photons absorbed by the
specimen. i.e.,
Particle flux J n
G= =
Ld Ld
1/ 2
J n0 meV Ê L ˆ meV [using Eqs. (ii) and (iv)]
Jn = = =Á ˜
e d Ë A¯ d
The multiplication and division of Ld in the above expression leads to
meV
Jn = (Ld )
( LA)1/ 2 d 2
Jn meV
\ G= = (viii)
Ld ( LA)1/ 2 d 2
The gain factor is also expressed in terms of transit time (Td) of an electron between the electrodes and life
time Te of the electron before recombination. Then the gain factor is defined as
G = Te/Td (ix)
A comparison of Eq. (ix) with Eq. (viii) yields
Te = (LA)–1/2
Td = d2 / meV
18.10 EffEcT Of TraPs

The trap is an energy level in the forbidden energy gap
of specimen which is capable of capturing either an
electron or a hole. The captured electron or hole may be
re-emitted at any time and can further move to another trap
(Fig. 18.14). There are two types of the traps which differ
from each other according to their operation. One of the
type helps recombining electron and holes for assisting to
achieve thermal equilibrium. The other type of trap does
not contribute directly in recombination process but they
contribute indirectly. Let us consider a specimen of the
crystal having N electron trap levels per unit volume at
low temperature. In the present aspect, Eq. (i) is
dn Figure 18.14
= L - An( n + N ) (x)
dt
where A is the recombination coefficient, which is the same for electron-hole recombination, and n is the
electron concentration in the conduction band. Thermal ionisation of carriers from traps back into the
conduction band contributes to the production rate through the term, say Bn . Hence Eq. (x) takes the form
1
dn (xi)
= L - An( n + N ) + Bn1
dt
In the steady state, we have
dn
= 0, n Æ n0 and Bn1 = 0. Therefore Eq. (xi) gives
dt
Ê Lˆ
n0 ( n0 + N ) = Á ˜
Ë A¯ (xii)
18.10.1 limiting cases
Now there are two limiting cases to be discussed. At low current, the carrier concentration n0 may be very
much less than electron trap levels (N), i.e., n0 << N.
Then from Eq. (xii), we have
L
n0 = (xiii)
NA
If n0 >> N, then from Eq. (xii) we get
1/ 2
Ê Lˆ
n0 = Á ˜
Ë A¯
18.10.2 response Time
The response time is related to the decay of carriers on switching off the light. This is the time in which the
carrier concentration falls to 1/e time of its initial value n0. This can be obtained by setting L = 0 in Eq. (x),
which yields
n+N n +N
ln - ln 0 = NAt (xiv)
n n0
Solution of Eq. (xiv) under the limit of N >> n0 is represented as
n = n0e–NAt (xv)
The time for the photocurrent to fall to e–1 of its initial value can be obtained from Eq. (xv) as
1
t0 = (xvi)
NA
A comparison of this equation with Eq. (vi) obtained for the case of absence of traps yields that the presence
of traps reduces the response time.
18.11 aPPlicaTiONs Of PHOTOcONDucTiviTy

There are many applications of photoconductivity, some of which are listed below.
(i) This is used to measure the intensity of illumination.
(ii) This is used to control the street light.
(iii) This is used in exposure meters for cameras.
(iv) This is used for relay control.
(v) This is used in moving object counters.
(vi) This is used in voltage regulator or systems.
(vii) This is also used for burglar alarms.
(viii) The important application of photocell, a device based on photoconducting, is in the reproduction of
sound in cinema, television and photography.
(ix) This is used in determination of opacity of solids and liquids.
(x) The temperature of celestial bodies is measured and their spectra are studied by using cells obtained
based on the concept of photoconductivity.
(xi) This is used to control temperature in furnaces and chemical reactions.
(xii) Photoelectric cells are used in obtaining electrical energy from sunlight during space travel. Thus,
the electric current is used to charge the battery.
s UmmarY

✦ When N atoms of the solid are well–separated, then these atoms lead to N-fold degenerate levels of the
solid. However, a continuously increasing interaction occurs between them when the atoms approach
one another to form a solid. This causes each of the levels to split into N distinct levels. Since a solid
contains about 1023 atoms per mole, i.e., N is very large, the splitted energy levels become so numerous
and close together that they form an almost continuous energy band. From the quantum point of view,
the wave functions of the valence electrons overlap and the overlapping of their wave functions results
in splitting or spreading of their energy levels.
✦ Free electron model ignores the effects those arise when the electrons interact with crystal lattice.
However, due to the effect of the periodic variation in the potential for an infinitely long lattice, the
wave function does not remain sinusoidal travelling waves of constant amplitude but now they include
the lattice periodicity in their amplitudes and electrons may be scattered by the lattice. In 1930, Kronig
and Penney proposed a one-dimensional model by assuming the shape of rectangular potential wells
and barriers along with the lattice periodicity. This is known as Kronig-Penney model.
✦ Kronig-Penney model tells about the discontinuities in energy when the wave number k satisfies the
condition k = np/a, where n = ±1, ±2, ±3,... etc. The E versus k diagram shows that an electron has
allowed energy values in the region between –p/a to +p/a. This region is called the first Brillouin zone.
Again there is another allowed energy zone, which is observed after the forbidden gap and is extended
from –p/a to –2p/a and p/a to 2p/a. This zone is called second Brillouin zone.
The concept of the Brillouin zones was extended to two-dimension by considering the wave number
k to have the two components kx and ky, i.e., along the x-axis and y-axis, respectively. This way kx =
ky =±2p/a limits the first Brilluion zone. The second Brillouin zone for a two-dimensional lattice is
represented by kx = ±2p/a and ky = ±2p/a excluding the space occupied by the first Brillouin zone.
✦ The concept of effective mass of an electron was introduced. This effective mass m* is given by
1 4p 2 d 2 E , where E is the energy of the electron and k is the wave number.
= 2
m * h dk 2
✦ During the formation of bands in solids it is observed that there is an energy gap, called forbidden band.
The forbidden gap represents energies which the electrons cannot occupy. Based on this energy gap and
the conduction, the solids are classified into different categories named as insulators, semiconductors
and conductors. A detailed description of these classes of the solids was given.
✦ The semiconductors in which the transformation of electrons to the conduction band and the generation
of holes in the valence band are achieved purely by thermal excitation are called intrinsic semiconductors.
During this effect equal numbers of electron and hole carriers are produced. The electrons and holes are
called intrinsic charge carriers and the resulting conductivity is known as intrinsic conductivity. Here
we obtained the electron concentration in the conduction band, hole concentration in valence band,
concentration of charge carriers, the energy band diagram and the Fermi level.
✦ The conductivity of an intrinsic semiconductor can be increased significantly by adding certain impurities
to it. By doing so, we get impurity semiconductor which is also known as extrinsic semiconductor. In
extrinsic semiconductors, current carriers are introduced by donor and acceptor impurities with locked
energy levels near the top or bottom of the forbidden gap.
The extrinsic semiconductors are of two types, namely n-type and p-type semiconductors. The concept
of Fermi energy was discussed in terms of energy diagram for both these types of the semiconductors.
✦ The effect of temperature was discussed on both n-type and p-type semiconductors. Since all the donors
donate their free electrons at room temperature, the additional thermal energy increases the generation
of electron-hole pairs. Thus the concentration of minority charge carriers increases. A temperature
is ultimately reached when the number of covalent bonds broken is very large such that the number
of holes and electrons is almost equal. The extrinsic semiconductor then behaves like an intrinsic
semiconductor, although its conductivity is higher. Thus with an increase in the temperature of an
extrinsic semiconductor, it behaves almost intrinsically.
✦ If a current carrying conductor is placed in a transverse magnetic field, a potential is developed in the
conductor in the direction perpendicular to both the current and magnetic field. This phenomenon is
known as Hall effect.
✦ Under Hall effect, with the application of magnetic field B, the accumulation of charge on the surfaces
of the specimen continues until the force on moving charges due to electric field associated with the
accumulated charge itself is large enough to cancel the force exerted by the magnetic field. So a steady
state condition is achieved. In this state, the value of electric field is called Hall electric field and the
JBd
potential difference so developed is called Hall voltage. The Hall voltage VH is given by VH = - ,
ne
where J is the current density, d is the width of the specimen (strip), n is the number of electrons per unit
1
volume and e is the electronic charge. The coefficient - is called the Hall coefficient RH.
ne
✦ Photoconductivity is an electrical phenomenon in which a material becomes more conductive due to
the absorption of light radiation. When the energy of the incident radiation is higher than the energy
gap Eg between conduction band and valence band, the electron–hole pairs are produced in the crystal.
The electrons are in the conduction band and the holes are in the valence band of the crystal. These
electron–hole pairs are the carriers of electrical conductivity.
✦ Gudden and Pohl discovered some basic experimental facts of the phenomenon of photoconductivity.
These are listed as follows.
(i) For a given material, the absorption of light and the excitation of photoconductivity by the light have
a similar dependence on the wavelength of light.
(ii) The region of photoresistivity gets extended to longer wavelengths in the presence of impurities.
✦ Simple model of photoconductor was discussed in detail by considering that the electron–hole pairs
are produced throughout the volume of the crystal when the light radiation falls on the crystal. In this
context, the response time t0 was defined as the time in which the electron concentration is dropped to
half of the number of electrons in the steady state, when the light falling on the specimen is stopped.
s
The response time t0 is given by t0 = , where s is the conductivity, me is the mobility of the
em e L
electrons and L is the number of photons absorbed by the crystal.
✦ We defined the sensitivity or the gain factor as the ratio of carriers crossing the specimen to the number
of photons absorbed by the specimen. The gain factor is also expressed in terms of transit time Td of
an electron between the electrodes and the life time Te of the electron before recombination. The gain
factor G is equal to Te/Td.
✦ The presence of impurities produces discrete energy levels in the forbidden gap. They are known as
traps. It means the trap is an energy level in the forbidden energy gap of specimen, which is capable of
capturing either an electron or a hole. The captured electron or hole may be re-emitted at any time and
can further move to another trap.
s olVeD e XamPles
E xamplE 1 Consider two-dimensional square lattice of side 3.0 Å. At what electron momentum values do the
sides of first Brillouin zone appear? What is the energy of free electron with this momentum?
Solution Given a = 3.0 ¥ 10–10 m.
Formula used for momentum of electron
p = k
p
For first Brillouin zone k = ± , then
a
h p h 6.62 ¥ 10 -34
p= = =
2p a 2a 2 ¥ 3 ¥ 10 -10
= 1.1 ¥ 10 -24 kg m/sec
p 2 (1.1 ¥ 10-24 ) 2
Energy E = =
2m 2 ¥ 9.1 ¥ 10-31
= 6.648 ¥ 10-19 J
= 4.155 eV
E xamplE 2 Find the position of Fermi level EF at room temperature (= 27oC) for germanium crystal having
5 ¥ 1022 atoms/m3.
Solution Given T = 27oC = 300 K and nC = 5 ¥ 1022 per m3
Formula used is
3/ 2
Ê 2p mkT ˆ
nC = 2 Á e( EF - EC ) / kT
Ë h 2 ˜¯
nC
e( EF - EC ) / kT = 3/ 2
Ê 2p mkT ˆ
2Á
Ë h 2 ˜¯
5 ¥ 1022
e( EF - EC ) / kT = 3/ 2
È 2 ¥ 3.14 ¥ 9.1 ¥ 10-31 ¥ 1.381 ¥ 10 -23 ¥ 300 ˘
2Í ˙
Î (6.62 ¥ 10-34 )2 ˚
5 ¥ 1022
=
25.115 ¥ 1024
e - ( Ec - EF ) / kT = 0.1991 ¥ 10 -2
e( EC - EF ) / kT = 502.296 or EC - EF = ln 502.296
EC - EF
= 6.2192 or EC - EF = 0.161 eV
kT
E xamplE 3 Consider the Fermi 0.3 eV below the conduction band at room temperature (=27oC) in an n-type
semiconductor. If the temperature is raised to 57oC, what would be the new position of Fermi level?
Solution Given EC – EF = 0.3 eV, T1 = 27oC = 300 K and T2 = 57oC = 330 K.
Formula used is
ÊN ˆ
EF = EC - kT ln Á c ˜
Ë Nd ¯
ÊN ˆ
EC - EF = kT ln Á c ˜
Ë Nd ¯
At temperature T1, EC - EF = kT1 ln Ê c ˆ

N
ÁË N ˜¯
d
At temperature T2, EC - EF¢ = kT2 ln Ê c ˆ

N
ÁË N ˜¯
d
EC - EF¢ T2
=
EC - EF T1
T2
or EC - EF¢ = ( EC - EF )
T1
330
= ¥ 0.3 eV
300
= 0.33 eV
E xamplE 4 For an intrinsic semiconductor having band gap Eg = 0.7 eV, calculate the density of holes and
electrons at room temperature (= 27oC).
Solution Given Eg = 0.7 eV.

In intrinsic semiconductor, the concentration of electrons and holes are same. So
3/ 2
È 2p kTm ˘
ne = nh = 2 Í e( EF - EC ) / kT
Î h 2 ˙˚
The Fermi level lies exactly in the middle of condition and valence band.
EC + EV
i.e. EF =
2
E + EV ( E - EV ) - E g
\ EF - EC = C - EC = - C =
2 2 2
3/2
È 2p kTm ˘ - E g / 2 kT
\ ne = nh = 2 Í e
Î h 2 ˙˚
3/2 È 0.7 ˘
È 2 ¥ 3.14 ¥ 1.38 ¥ 10-23 ¥ 300 ¥ 9.1 ¥ 10-31 ˘ -Í ˙
Î 2 ¥ 0.026 ˚
=2¥Í ˙ e
Î 6.662 ¥ 10-34 ˚
= 3.6 ¥ 1019 per m 3
E xamplE 5 Assuming that there are 5 ¥ 1028 atoms/m3 in copper, find the Hall coefficient.
Solution Given n = 5 ¥ 1028 atoms /m3.
Formula used is
1 1
RH = - =
ne 5 ¥ 1028 ¥ 1.6 ¥ 10 -19
= - 0.125 ¥ 10 -9 m 3 /C
E xamplE 6 Using free electron model, find the Hall coefficient of sodium assuming bcc structure for Na of
cell side 4.28 Å.
Solution Given a = 42.8 ¥ 10–10 m.
Unit cell of sodium atom (Na) of volume a3 has 2 atoms, i.e.,
1 2
n=2 3 = = 2.55 ¥ 1028 per m3
a (4.28 ¥ 10 -10 )3
1 -1
Hall coefficient RH = - =
ne 2.551 ¥ 10 ¥ 1.6 ¥ 10 -19
28
= - 0.245 ¥ 10 -9 m 3 /C
Q.1 The energy eigen value in a free electron model is given by

h2 h2 k 2 2 k 2 2 2
(a) E = 2 2
(b) E = (c) E = (d) E = 2m k
8p k m 2m 2m
Q.2 The first Brillouin zone is defined between the region
p 2p p
(a) k = 0 to (b) - to
a a a
p 2p p p
(c) k = - to (d) - to
a a a a
Q.3 The second Brillouin zone is defined between the region
-p +p -2p -p p 2p
(a) k = to (b) k = to and to
a a a a a a
-p 2p -2p 2p
(c) k = to (d) k = to
a a a a
Q.4 Bragg’s conditions for diffraction is given by the relation

(a) 2 sin q = nl (b) 2a sin q = l
(c) 2a sin q = nl (d) a sin q = nl
Q.5 Which one of the following relations is correct for effective mass of an electron
-1
h2 d 2 E * h 2 È dk 2 ˘
*
(a) m = (b) m = Í ˙
4p 2 dk 2 4p 2 Î d 2 E ˚
-1 -1
* È d 2E ˘
2 h2 È d 2 E ˘
(c) m = h Í 2 ˙
*
(d) m = Í ˙
Î dk ˚ 4p 2 Î dk 2 ˚
Q.6 When we increase the temperature of extrinsic semiconductor, after a certain temperature it behaves
like
(a) an insulator (b) an intrinsic semiconductor
(c) a conductor (d) a superconductor
Q.7 Pure semiconductor behaves as an insulator at
(a) 273 K (b) –273oC (c) 373K (d) none of these
Q.8 The group velocity of the electron in one–dimensional lattice is defined as
dk dw w
(a) vg = (b) vg = (c) vg = (d) vg = wk
dw dk k
Q.9 Hall coefficient is given by the relation
J 1 -J
(a) RH = –neJ (b) RH = (c) RH = - (d) RH = -
ne ne ne
Q.10 Which type of semiconductor material has negative Hall coefficient
(a) p-type (b) n-type (c) intrinsic (d) none of these
Q.11 The electrical conductivity of pure semiconductor at 0K is
(a) zero (b) large (c) finite (d) none of these
Q.12 The electrical current in an intrinsic semiconductor is because of
(a) electrons in conduction band (b) holes in valance band
Q.13 The majority carriers in n-type semiconductors are
(a) electrons in conduction band (b) holes in valence band
Q.14 By doping fifth group elements in germanium crystal, which type of semiconductor material is
obtained
(a) positive (b) negative (c) neutral (d) none of these
Q.15 In an n-type semiconductor the Fermi-level lies
(a) above the top of the valence band (b) below the bottom of the conduction band
(c) In the middle of the forbidden gap (d) near the conduction band
Q.16 On increase of temperature, the Fermi level shifts upwards in
(a) p-type semiconductor (b) n-type semiconductor
(c) intrinsic semiconductor (d) none of these
Q.17 The concentration of electrons in the conduction band of an intrinsic semiconductor is proportional to
(a) T (b) T2 (c) T3/2 (d) T3
Q.18 The energy gap between the valence and conduction bands in a semiconductor is of the order of
(a) 26 eV (b) 1.0 eV (c) 7.0 eV (d) 0.001 eV
Q.19 Which of the following relation is correct for product of electron-hole concentration (np)
(a) AT 3e–Eg/KT (b) AT 5e–Eg/KT (c) AT 7e–Eg/KT (d) none of these
Q.20 The product of electron-hole concentration (np) changes with
(a) temperature (b) doping concentration
(c) pressure (d) none of these
Q.21 Which one of the following is correct for Kronig-Penney model
(a) real model (b) approximate model
(c) virtual model (d) none of these
Q.22 Which of the following theory proposes about the free electron inside periodic lattice
(a) zone theory (b) quantum theory of free electron
(c) classical theory of free electron (d) none of these
Q.23 In the absence of potential, per Kronig-Penney model
(a) forbidden energy regions are not there (b) forbidden energy regions are there
(c) all values of energy are allowed (d) none of these
Q.24 In the absence of potential barrier, the E-k curve is
(a) parabola with discontinuities (b) continuous parabola
(c) discontinuous energy levels (d) none of these
Q.25 The discontinuities occur in E-k curve at
np np a
(a) k = ± 2 (b) k = ± (c) k = ± (d) none of these
a a np
Q.26 The effective mass of an electron may be
(a) positive (b) negative (c) infinity (d) all of these
Q.27 The Hall-effect is used in determining
(a) mobility of charged carriers (b) density of charged carriers in extrinsic semiconductor
(c) type of extrinsic semiconductor (d) all of these
Q.28 In the phenomenon of photoconductivity a material becomes
(a) more conductive (b) more resistive
(c) less conductive (d) none of these
Q.29 The phenomenon of decrease of resistivity of an insulator when light radiation falls on it, is called
(a) photoelectric effect (b) Frenkel effect
(c) Compton effect (d) photoconductivity
Q.30 The photoconductivity is mainly due to
(a) extrinsic excitations (b) intrinsic excitations
Q.31 The minimum energy which is required for intra excitation is
(a) greater than the forbidden gap Eg (b) less than the forbidden gap Eg
(c) equal to the forbidden gap Eg (d) none of these
Q.32 Which one of the following statements is not correct for photoconductive cells?
(a) they are made of a single photoresistive material
(b) they are also called photoresistors
(c) they have a forward biased p–n junction
(d) they have a high dark-to-light resistance ratio
Q.33 The photocurrent density of a photoconductor is given by the relation
d n0 eme n em V
(a) J + (b) J (c) J = 0 e (d) J = n0emeVd
n0 emeV dV d
Q.34 The presence of traps in a photoconductor has the following effect
(a) it reduces the response time (b) it increases the response time
(c) it has no effect on response time (d) it has no relevance with photoconductor
Q.35 Which one of the following relations is correct for response time of a photoconductor
s 1 sm e
(a) t0 = semeL (b) tq = (c) tq = (d) t0 =
em e L s eme L eL
Q.36 In a photoconductive cell, the internal resistance changes with a change in
(a) frequency of light (b) intensity of light
(c) both of these (d) none of these
Q.37 Photoconductor is also known as
(a) photovoltaic cell (b) solar cell (c) photoresistor (d) photodiode
P ractice P roblems
General Questions
Q.1 Explain how the atomic energy levels split into bands when a number of atoms are brought close
together to form a crystal?
Q.2 Discuss Kronig-Penney model. Using the model show the energy spectrum of electron consisting of a
number of allowed energy bands separated by forbidden bands.
Q.3 What is the effect of periodic potential on the energy of electrons in a metal? Explain it on the basis of
Kronig-Penney model and explain the formation of energy bands.
Q.4 Discuss how the concept of bands was originated in solids. Give necessary theory. What is E-K diagram
and what do you infer from them?
Q.5 What are Brillouin zones? Explain using E-K diagrams.
Q.6 Discuss the formation of Brillouin zones for (i) linear lattice (ii) two–dimensional lattice.
2
Q.7 Define (m*) and prove that effective mass of an electron m* = 2
d E / dk 2
Give the physical basis of effective mass and explain its physical significance.
Q.8 Write notes on following
(a) Intrinsic and extrinsic semiconductors
(b) Effective mass
Q.9 Explain origin of bands in solids.
Q.10 Write short note on band theory of solids.

Q.11 Discuss the band theory of solids and explain the formation of bands and concept of holes.
Q.12 Discuss the motion of electrons in a periodic field of a crystal and show that the effective mass of
an electron in a crystal is inversely proportional to the second derivative of E-K curve. Under what
conditions, the effective mass of an electron can become +ve, –ve and infinity?
Q.13 Based on band theory of solids, distinguish between conductors, semiconductors and insulators.
Q.14 Name three semiconductors along with the value of band gaps.
Q.15 Derive an expression for the carrier concentration in an extrinsic semiconductor. What would be the
position of Fermi level? Explain.
Q.16 Write short notes on intrinsic and extrinsic semiconductors.
Q.17 Show that Hall coefficient is independent of the applied magnetic field and is inversely proportional to
the current density and electronic charge.
Q.18 What is Hall effect? Give an experimental method of calculating concentration and type of charges in
a given semiconductor.
Q.19 Write short note on Hall effect and its application.
Q.20 Explain how Hall effect helps in determining the sign of charge carrier in a material. Also prove that
f = mH where f is the Hall angle, m is mobility of charge carrier and H is the applied magnetic field.
Q.21 What is photoconductivity?
Q.22 State the principle of photoconductive cell. Describe its construction, working and uses. Show the
illumination characteristics and spectral response.
Q.23 What are traps? Discuss the effect of traps on photoconductivity of a photoconductor.
Q.24 What are traps? Discuss a simple model to show the effect of traps on the photoconductivity.
Q.25 Why does the electrical conductivity increase when certain solids are exposed to light of suitable
wavelengths? Suggest simple model of a photoconductor and explain the following (a) gain,
(b) response time, (c) effect of traps.
Q.26 Explain with proper diagram the working of a simple model of photoconductor.
Q.27 Discuss sensitivity or gain factor in detail.
Magnetic Properties
of Solids
19
Learning Objectives
LO 1 Define magnetic susceptibility, relative LO 5 Know about classical theory of
permeability and magnetic moment of ferromagnetism
an electron LO 6 Discuss hysteresis, energy loss due to
LO 2 Discuss classification of magnetic hysteresis, and importance of hysteresis
materials curve
LO 3 Compare properties of paramagnetic, LO 7 Learn about magnetic circuits
diamagnetic and ferromagnetic LO 5 Analyse magnetic materials and their
materials applications
LO 4 Explain classical theory of
diamagnetism and paramagnetism
Introduction
All materials, i.e., metals, semiconductors and insulators, reveal the phenomenon of magnetism.
Magnetic materials play an important role in modern technology as they are frequently used in industrial
electronics, computer industry, etc. The traditional methods of information storage and retrieval are
rapidly being replaced by magnetic storage. The magnetism of materials is mainly an outcome of the
interactions of magnetic moments of their constituent atoms or molecules. Magnetic materials can be
classified into three categories on the basis of their permeability or susceptibility. The magnetic materials
for which susceptibility cm is negative (permeability mr £ 1) are said to be diamagnetic materials, whereas
the materials with positive susceptibility (permeability mr ≥ 1) are said to be paramagnetic materials. If the
susceptibility is much larger than zero and permeability mr >> 1, then the magnetic materials are called
Æ
ferromagnetic materials. Paramagnetic and diamagnetic materials have a linear relationship between B
Æ Æ Æ
and H . However, the relationship between B and H is nonlinear for ferromagnetic materials. A material
is said to be nonmagnetic if susceptibility cm = 0 (or mr = 1), it is magnetic otherwise. Depending on the
alignment of magnetic moments within the materials, these are further classified into five important groups,
namely, diamagnetic, paramagnetic, ferromagnetic, anti-ferromagnetic and ferrimagnetic materials. Since
diamagnetism, paramagnetism and antiferromagnetism are weak effects, the materials which exhibit
these phenomena are known to be nonmagnetic. However, ferromagnetism and ferrrimagnetism are very
strong effects. Therefore, in a large number of devices, these two magnetic phenomena are prominently
utilised.
The magnetic materials are of two types, namely, soft materials and hard materials. Soft magnetic materials
are used in ac applications, since they are easily magnetised and demagnetised. However, hard magnetic
materials are used in producing permanent magnets, since they retain magnetism on a permanent basis.
Due to such properties, these materials are significantly used in information storage devices. In order to
realise the operating principles of different magnetic devices, it is essential to understand the magnetic
phenomena. So at first we define various terms, viz., intensity of magnetisation, magnetic susceptibility,
Æ Æ
relative permeability, etc. Magnetic flux density B and magnetic field strength H have already been
discussed in detail in Chapter 10.
Intensity of Magnetisation (I)

It is defined as the magnetic moment per unit volume of the magnetised substance
M
I=
V
which for a substance of length 2l and cross-sectional area a becomes
m ¥ 2l m
I= =
a ¥ 2l a
Thus, it can also be defined as pole-strength per unit area of cross-section. The intensity of magnetisation is
sometimes represented by M. In that case, another symbol is used for the magnetic moment.
Magnetic Susceptibility (cm)

It is the ratio of the magnetic moment per unit volume (I) to the magnetic field strength (H) of the magnetising
field.
I
cm =
H
It is positive for a paramagnetic material and negative for a diamagnetic one.
Relative Permeability (mr)

It is the ratio of the magnetic permeability (m) of the substance to the permeability of the free space (m0).
m
mr =
m0
This can also be defined as the ratio of the magnetic flux density produced in the medium to that which would
be produced in a vacuum by the same magnetising force.
Magnetic Properties of Solids 687
Relation between Permeability (mr) and Magnetic Susceptibility (cm)

As discussed earlier, the magnetic flux density B can be written in terms of the magnetic field strength H and
the intensity of magnetisation I as
B = m0(H + I) (i)
I
cm = (ii)
H
Therefore,
Ê Iˆ
B = m0 ( H + I ) = m0 H Á1 + ˜ = m0 H (1 + c m )
Ë H¯
B
or = m = m0 (1 + c m )
H
m
\ Relative permeability m r = = 1 + cm
m0
where m0 is the magnetic permeability of free space.
19.1 MagnetIc MOMent Of an eLectROn LO1

The magnetic properties of solids originate due to the motion of electrons. An atom is made up of a positively
charged nucleus placed at the centre and negatively charged electrons which move around the nucleus. The
orbital motion of each electron around the nucleus may be treated as a current loop. As a result, it sets up a
magnetic field.
The current produced,
charge e
i= = (i)
time period T
distance 2p r
T= = (ii)
velocity v
where v is the linear velocity of the electron in a circular orbit of radius r. Then,
e ev
i= = (iii)
2p r /v 2p r
Magnetic moment of a magnetic dipole is
ev 1
M = iA = p r 2 = evr (iv)
2p r 2
According to Bohr, the magnetic moment can also be expressed as
nh
mvr =
2p
nh
or vr = (v)
2p m
By putting the value of vr from Eq. (v) in Eq. (iv), we get
e Ê nh ˆ È eh ˘
M = Á ˜ = nÍ , where n = 1, 2, 3, …
2 Ë 2p m ¯ Î 4p m ˙˚
eh
The above relation gives the magnetic moment of an electron orbiting around a nucleus. The quantity
4p m
is called Bohr magneton, represented by mB.
19.2 cLaSSIfIcatIOn Of MagnetIc MateRIaLS LO2

Magnetic materials can be classified into five different categories, namely, diamagnetic materials, paramag-
netic materials, ferromagentic materials, anti-ferromagnetic materials and ferrimagnetic materials. We dis-
cuss below all the materials in detail.
19.2.1 Diamagnetic Materials
On placing in an external magnetic field, the materials which acquire feeble magnetism in the direction opposite
to that of the applied field are called diamagnetic materials. This property is found in the substances whose
outermost orbit has an even number of electrons. Since the electrons have spins opposite to each other, the net
magnetic moment of each atom is zero. The magnetism of diamagnetic materials is called diamagnetism. If these
materials are brought close to the pole of a powerful electromagnet, they are repelled away from the magnet.
Examples of diamagnetic materials are bismuth, zinc, copper, silver, gold, lead, water, etc.
19.2.2 Paramagnetic Materials
On placing in an external magnetic field, the materials which acquire feeble magnetism in the direction of
the applied field are called paramagnetic materials, and their magnetism is known as paramagnetism. This
property is found in the substances whose outermost orbit has an odd number of electrons. The source of
paramagnetism is the permanent magnetic moment possessed by the atoms of paramagnetic materials. If these
substances are brought close to a pole of a powerful electromagnet, they get attracted towards the magnet.
Examples of paramagnetic materials are aluminium, odium, platinum, manganese, copper chloride, liquid
oxygen, etc.
19.2.3 ferromagnetic Materials
On placing in an external magnetic field, the materials which acquire strong magnetism in the direction of the
applied field are called ferromagnetic materials and their magnetism is called ferromagnetism. This property
is found in the substances which are generally like paramagnetic materials. These are strongly attracted by
magnets.
Examples of ferromagnetic materials are iron, nickel, cobalt, magnetite (Fe3O4), etc.
19.2.4 anti-ferromagnetic Materials
Anti-ferromagnetic substances are crystalline materials. In these materials, the dipole moments of the
neighbouring dipoles are equal and opposite in orientation so that the net magnetisation vanishes. If they
are placed in the magnetic field, they are feebly magnetised in the direction of the field. Such materials
are called anti-ferromagnetic materials and their magnetism is called anti-ferromagnetism. Examples
of anti-ferromagnetic materials are: MnO, FeO, CaO, NiO, MnO4, MnS, etc. Susceptibility of these
materials vary with temperature. It increases with increasing temperature and reaches a maximum at a
particular temperature called the Neel temperature (TN). Above this temperature, these materials behave like
paramagnetic materials.
19.2.5 ferrimagnetic Materials
If the spins of the atoms are such that there is a net magnetic moment in one direction, the materials are called
ferrimagnetic materials. Examples of ferrimagnetic materials are ferrites which consist of mainly ferric oxide
Fe2O3 combined with one or more oxides of divalent metals.
19.3 cOMPaRISOn Of PROPeRtIeS Of PaRaMagnetIc, DIaMagnetIc LO3

anD feRROMagnetIc MateRIaLS
S.No. Paramagnetic Materials Diamagnetic Materials Ferromagnetic Materials

1. These materials show positive These materials show negative These materials show positive
magnetic susceptibility (of the susceptibility (of the order of and high magnetic susceptibility
order of 10–6). 10–6). (~106).
2. The relative permeability is slightly mr is slightly less than unity The mr for a ferromagnetic
more than unity (mr > 1). (mr < 1). material is of the order of few
thousands.
3. The magnetic susceptibility The magnetic susceptibility The magnetic susceptibility
is strongly dependent on of diamagnetic materials is decreases with increase in
temperature and varies inversely practically independent of temperature.
with temperature. temperature.
4. When a bar of a paramagnetic When a bar of these materials is When a bar of these materials
material is suspended between the suspended between the poles of is suspended between the poles
poles of a magnet, it stays parallel a magnet, it stays parallel to the of a magnet, it behaves like a
to the lines of force. magnetic field. paramagnetic material.
5. If these materials are placed in If these materials are placed in These materials behave like
a non-uniform field, they are a non-uniform field, they are paramagnetic substances, if
attracted towards the stronger attracted towards the weaker field. placed in a non-uniform field.
field.
19.4 ClassiCal Theory of DiamagneTism (langevin’s Theory) LO4

The classical theory of diamagnetism was developed by the French Physicists Paul Langevin in 1905. We
have discussed earlier that atoms of diamagnetic materials do not possess a permanent magnetic moment.
When the material is placed in an external magnetic field, the orbital motion of electrons changes and hence
atoms acquire an induced magnetic moment. The direction of the induced dipole will be opposite to that of
the external magnetic field.
Suppose in the absence of an external magnetic field, an electron of mass m moves with linear
velocity v0 (or angular velocity w0) in a circular orbit of radius r (Fig. 19.1).
r
The magnetic moment equivalent to this motion is given by
M = iA (i)
Figure 19.1
By the definition of current,

e e
i= = = ne, where n is the frequency
T 1/n
w e
or i= 0 (ii) [\ w0 = 2pn]
2p
with the help of Eqs. (i) and (ii), we have
ew 0 ew
M = A = 0 p r2
2p 2p
er 2w 0
or M = (iii)
2
Under the above condition (when B = 0), the centripetal force on the electron is given by
mv02
Fc = = mw 02 r (iv)
r
Now, if an external magnetic field B is applied to the plane of the electron orbit normally Fig. 19.2, an
additional magnetic force will act on the electron, given by
Æ Æ Æ
F m = –e(v ¥ B ) (v)
fi Fm = –evB
Fm
× × × × × × × ×V × × × ×
× × ×Fm × × × × × × × × ×
× × × F×c × × × × × F×c × ×
× × × × × × × × × × × ×
× × × × × × × × × × × ×
(a) (b)
Figure 19.2
By Fleming’s left hand rule, this magnetic force acts on the electron radially inward or outward if the electron
moves in clockwise Fig. 19.2a or anti-clockwise Fig. 19.2b direction, respectively. Hence, the total force on
the electron will be
Æ Æ Æ
Ft = Fc ± Fm
Æ
F t = mw02r  evB (vi)
Due to the magnetic force, the angular frequency changes from w0 to w. This change in angular frequency can
be calculated as under. By Faraday’s law of induction
df
|e | = - (vii)
dt
According to the definition of induced e.m.f.
 
e= Ú E ◊ dl = E 2p r
e
or E= (viii)
2p r
where r is the radius of the orbit.
By Newton’s second law of motion
dv
F=m = - eE (ix)
dt
By using Eqs. (viii) and (ix), we get
dv e
m =- e (x)
dt 2p r
By using Eqs. (vii) and (x), we get
dv e df
m =-
dt 2p r dt
e
or dv = - df
2p rm
e
or Dv = - Df
2p rm
As the external magnetic field changes from 0 to B, the corresponding flux changes from 0 to pr2B.
Then
e e
Dv = - (p r 2 B - 0) = - p r2B
2p rm 2p rm
eBr
or Dv = - (xi)
2m
Therefore, the change in angular velocity will be
Dv eB
Dw = =- (xii)
r 2m
This change in angular frequency is also known as Larmour frequency.
On application of an external magnetic field, the angular frequency of the orbital electron gets changed
which leads to a charge in magnetic moment. Therefore, in both the cases of electron rotating in clockwise
or anti-clockwise directions, the magnetic moment of the orbital electron changes and hence it can be
obtained as
er 2
DM = Dw [Using Eq. (iii)]
2
er 2 È eB ˘
= -
2 ÍÎ 2m ˙˚
e2 B 2
\ DM = - r (xiii)
4m
As seen above, the negative sign shows that the induced magnetic moment is always opposite to change in the
magnetic field. In deriving the above equation, we have assumed that the orbit of the electron is normal to the
applied field. But these orbits can have any orientation with the field. Since there are a number of randomly
oriented electron’s orbits in the atom that show the spherical symmetry, the total induced magnetic moment
in the atom is given by
Ze 2 B 2
DM = - r
4m
Æ
Here Z is the atomic number (i.e., the number of electrons in an atom) and r 2 is the mean square radius of the
electrons orbits. Let (x, y, z) be the coordinates of any point on the spherical orbit of radius r. Then
r2 = x2 + y2 + z2
Again consider x , y , z– as the average values of components of radii for all the electrons along the three axes. Then
– –
r 2 = R2 = x–2 + y–2 + z–2

Æ
For the spherically symmetric charge distribution,

R2
x2 = y2 = z 2 =
3
If the direction of the magnetic field is along the z-axis then considering the plane normal to the direction of
magnetic field we get
2R2
r 2 = x2 = y2 =
3
Therefore, the total induced magnetic moment in the atom will be
Ze 2 B 2 R 2
DM = - ¥ (xiv)
4m 3
Ze 2 R 2 B
\ DM = -
6m
If N is the number of atoms per unit volume of the substance, then the induced magnetic moment per unit
volume is
NZe 2 R 2 B
I induced = - (xv)
6m
Since B = m0H,
NZe 2 R 2 m0 H
I induced = -
6m
By definition, the magnetic susceptibility of a diamagnetic substance is
I induced m NZe 2 R 2
cm = =- 0 (xvi)
H 6m
From the above equation, it is clear that the diamagnetic susceptibility is negative, which does not depend on
the temperature and the intensity of the external magnetic field.
19.5 ClassiCal Theory of ParamagneTism (langevin’s Theory) LO4

The classical theory of paramagnetism was developed by Langevin in 1905. In paramagnetic materials, each
individual atom has a permanent magnetic moment. This is because of the fact that the magnetic moment of the
individual electrons due to their orbital as well as spin motions do not cancel out. So an individual atom has a
permanent dipole. This property is found in those substances whose outermost orbit has an odd number of
electrons. In the absence of an external magnetic field, these atomic magnets are randomly oriented in the
substance and the total magnetic moment of the substance is zero. When a magnetic field
B
is applied, these atomic dipoles tend to line up in the direction of the external magnetic
M
field but the temperature opposes this alignment. q
e
Let the paramagnetic substance contain N atoms per unit volume at a temperature T(K).
Each atom possesses a permanent magnetic moment M. The potential energy of an
atomic magnetic dipole inclined at an angle q with direction of external field, as shown
Figure 19.3
in Fig. 19.3 is
U = −MB cos q (i)

According to Boltzmann’s classical law, the number of atoms per unit volume in the energy ranging from U
to U + dU is given by
dN = Ce–U/k T dU
= CeMB cos q/k T MB sin q dq (ii)
where k is the Boltzmann constant and C is the constant of proportionality depending on the atoms. Integrating
Eq. (ii) for q from 0 to p, we can find total number of atoms per unit volume of the substance as
p
N = Ú dN
0
p
= C Ú e MB cosq / kT MB sin q dq (iii)
0
Thus,
N (iv)
C= p
Úe
MB cos q / kT
MB sin q dq
0
We know that M is the magnetic moment of each magnetic dipole that makes an angle q with the direction of
the external magnetic field. So, its component in the direction of external field will be M cos q. Thus, resultant
magnetic moment due to atoms along the external field will be M cos q dN. Hence the total magnetic moment
per unit volume of the substance (i.e., the intensity of magnetisation) is given by
p
I = Ú M cos q dN (v)
0
By using Eqs. (ii) and (v), we get

p
I = Ú M cos q Ce MB cosq / kT MB sin q dq
0
p
I = CM 2 B Ú e MB cosq / kT sin q cos q dq (vi)
0
By putting the value of C in Eq. (vi) from Eq. (iv), we have

p
Úe
MB cos q / kT
sin q cos q dq
2
I = NM B 0
p
Úe
MB cos q / kT
MB sin q dq
0
p
Úe
MB cos q / kT
sin q cos q dq
= NM 0 (vii)
p
Úe
MB cos q / kT
sin q dq
0
MB
By putting = x, cos q = y, so –sin q dq = dy in Eq. (vii), we have
kT
-1 -1 +1
- Ú e xy y dy Ú e y dy
xy
Úe
xy
y dy
+1 +1 -1
I = NM -1
= NM -1
= NM +1
- Ú e dy xy
Úe
xy
dy Úe
xy
dy
+1 +1 -1
+1
È ye xy e xy ˘
Í - 2˙
Î x x ˚ -1 È (e x + e - x ) 1 ˘
I = NM +1
= NM Í x -x
- ˙
È e xy ˘ Î (e - e ) x ˚
Í ˙
Î x ˚ -1
È 1˘
or I = I 0 Ícot hx - ˙ (viii)
Î x˚
where NM = I0 shows the saturation value of the intensity of
1.0
magnetisation (I) when all the magnetic dipoles get aligned in the
direction of the external magnetic field.
I/I0
In Eq. (viii), the function I / I0 = cot hx –1/x is called the Langevin
0.5
function and is represented by L(x). The variation of L(x) with x is
shown in Fig. 19.4. From the figure, it clear that
Ê MB ˆ
(i) If x Á = is large, i.e., temperature is very low, then
Ë kT ˜¯ x
Figure 19.4
L( x) ª cot hx
or I = I0
Thus, for a very low temperature (or in the strong magnetic field) all the magnetic dipoles are
aligned in the direction of external field and a saturation is obtained.
Ê MB ˆ
(ii) If x Á = is very small i.e., the temperature is very high or in other words, the external magnetic
Ë kT ˜¯
field is weak, then we have
È 1˘ È1 Ê x2 ˆ 1 ˘
I = I 0 Ícot hx - ˙ = I 0 Í Á1 + ˜ - ˙
Î x˚ Îx Ë 3 ¯ x˚
x MB MB
= I0 = I0 = NM
3 3kT 3kT
NM 2 NM 2
I= B= mH (ix)
3kT 3kT
I m NM 2 m ( NM ) 2
cm = = = [Q NM = I 0 ]
H 3kT 3 NkT
m I 02
or cm = (x)
3 NkT
Eq. (x) is known as the Curie law, which can also be written as
C
cm =
T
m I 02
where C = is known as Curie constant.
3 Nk
This equation shows that the magnetic susceptibility cm of a paramagnetic material depends on the temperature
T and it varies inversely with T.
19.5.1 failure of langevin’s Theory

This theory was not able to explain the complicated dependence of susceptibility on the temperature as shown
by several paramagnetic substances. In view of this, Langevin’s theory was further modified by Curie and
Weiss.
19.5.2 curie–Weiss Law
In 1907, Weiss had modified the Langevin’s theory of paramagnetism. He assumed that in a paramagnetic
substance an internal molecular magnetising field is generated because of mutual interaction between
the atomic magnetic dipoles. If the molecular magnetising field (Hi) generated at any point because of a
neighbouring atomic magnet is proportional to the intensity of magnetisation (I), then we have
Hi µ I or Hi = lI (xi)
where l is the molecular field coefficient and is independent of temperature. Hence, the effective magnetising
field within the substance may be expressed as
Heffective = H + lI (xii)
With which Eq. (ix) becomes
NM 2 m H effective NM 2 m ( H + l I )
I= =
3kT 3kT
Ê NM 2 ml ˆ NM 2 m H
or I Á1 - ˜=
Ë 3kT ¯ 3kT
NM 2 m H
or I= (xiii)
3kT - NM 2 ml
The magnetic susceptibility then becomes
I NM 2 m NM 2 m
cm = = =
H 3kT - NM 2 ml È NM 2 ml ˘
3k ÍT - ˙
Î 3k ˚
NM 2 m
3k C È m I 02 NM 2 m ˘
cm = = ÍQ C = = ˙ (xiv)
NM 2 ml T - q Î 3 Nk 3k ˚
T-
3k
NM 2 ml
where q=
3k
The relation (xiv) is called the Curie–Weiss law and the constant q is known as Curie temperature. It is clear
from the relation that if T < q, then the magnetic susceptibility of the paramagnetic substance becomes negative
and it behaves like a diamagnetic substance. Hence, Curie–Weiss law is applicable only for temperatures
T > q.
19.6 cLaSSIcaL theORy Of feRROMagnetISM LO5

We know that each atom of the ferromagnetic materials such as iron, nickel, cobalt, etc. has a permanent
magnetic moment like paramagnetic substances. The magnetic susceptibility of a ferromagnetic substance is
a thousand times greater than that of a paramagnetic substance. In general, a specimen of a ferromagnetic
substance contains a number of small regions called domains. These domains contain a large number of
atoms, nearly 1017 to 1022, and have the dimensions of about 10–6 cm3 to 10–2 cm3. Every domain is
magnetically saturated and the direction of magnetisation in different domains is different.
If the external magnetic field is absent, all the magnetic
domains in the substance are randomly oriented and their
resultant magnetic moment in any direction will be zero. As
per modern theory, it is assumed that in the absence of an
external magnetic field, these domains form closed loops
within the substance so that the net magnetic moment of the Figure 19.5
whole substance is zero, as shown in Fig. 19.5a.
When a magnetic field is applied to the ferromagnetic substance, the substance becomes magnetised. The
intensity of magnetisation of the substance depends upon the intensity of an external magnetic field.
When the substance is placed in a weak external magnetic field, the Imax
magnetisation produced is due to the displacement of boundaries of domains Q R
Fig. 19.5b and if the external magnetic field is strong, the magnetisation
I
produced is mainly by the rotation of domains Fig. 19.5c. Fig. 19.6 represents
the magnetisation curve for the ferromagnetic substance. In a very weak
magnetic field, as represented in the part OP of the curve, the displacement of P
boundaries of domains is reversible and if we removed the external magnetic
field, the boundaries of domains again come back to their original positions. 0
B
If we increase the external magnetic field, as represented in the part PQ of Figure 19.6
the curve, the displacement of boundaries of domains is irreversible and the
material immediately becomes magnetised. If we again increase the magnetic field, as represented in the
part QR of the curve, the magnetisation of the substance is because of rotation of domains in the direction of
magnetising field.
Thus, the net effective magnetic field of the ferromagnetic substance is given by
Heffective = H + Hi
where Hi is the magnetic field generated due to the mutual interaction between magnetic dipoles. By using
Curie–Weiss law as represented in Eq. (xiv), the magnetic susceptibility of the ferromagnetic substance is
C
cm =
T -q
where C and q are the Curie constant and the Curie temperature, respectively.
On the basis of the above relation, the following conclusions can be drawn.
(a) If T = q, the magnetic susceptibility will approach to infinity.
(b) If T < q, the magnetic susceptibility will be negative. In this condition, the Curie–Weiss law is not
applicable because the ferromagnetic substance gets magnetised even in the absence of external
magnetic field.
(c) If T > q, the magnetic susceptibility decreases with the increase in temperature. In this condition, the
ferromagnetic properties disappear and the substance becomes paramagnetic.
Æ Æ
19.7 hysTeresis: nonlinear relaTionshiP beTween b anD h LO6
Ferromagnetic materials like iron and steel are used for screening (or shielding) that protect sensitive electrical
devices from disturbances from strong magnetic fields. An example of an iron shield is the compass, which
without shielding gives an erroneous reading due to the effect of external magnetic field. For perfect screening,
it is required that the shield have infinite permeability (mr = •). B
Æ Æ Æ
Bmax P
Even though B = m0(H + I ) holds good for all materialsÆ Æ
including Permanent flux Virgin curve
ferromagnetic materials, the relationship between B and H depends density Br
(Retentivity)
on previous magnetisation of a ferromagnetic material, i.e., on its Æ −HC
magnetic
Æ
history.
Æ
Instead
Æ
of having a linear relationship between B HC Hmax H
O
and H (i.e., B = mH ), it is only possible to represent the relationship −Br
Coersive field intensity
(Coercivity)
by a magnetisation curve or a B–H curve, as shown in Fig. 19.7. At
−Bmax
any point on this curve, m is given by the ratio B–H, and not by the Q
slope of the curve (dB / dH). Figure 19.7

We can explain the B–H curve as follows. Initially, a ferromagnetic material is unmagnetised. As H is increased
due to increase in the current from O to the maximum applied field intensity Hmax, a curve OP is produced.
This curve is known as the virgin or initial magnetisation curve. Now we move back and decrease H. It is seen
that when H decreases after P, B does not follow the initial curve but lags behind H. This is called hysteresis.
When H reaches zero, it is obtained that B π 0, i.e., the material possesses some finite B. This finite Br is
called the permanent flux density or residual magnetism which depends on Hmax. The power of retaining
this magnetism is called the retentivity of the substance. It is a measure of the remaining magnetisation in
the substance when the magnetising field is removed. The existence of Br is the cause of having permanent
magnets. At H = Hc (decreased by reversing the current) B = 0. This value of H, i.e., Hc, is called the coercive
field intensity or coercivity of the substance. It is a measure of the reverse magnetising field required to
destroy the residual magnetism of the substance. Materials for which Hc is small are said to be magnetically
hard. Hc also depends on Hmax. Further increase in H to reach Q and in reverse direction to reach P gives
a closed hysteresis loop. The shape of this loop varies from one material to another. Some ferrites have an
almost rectangular hysteresis loop. These ferrites are used in digital computers as magnetic information
storage devices. The area of the loop represents energy loss (hysteresis loss)/ unit volume during one cycle of
the periodic magnetisation of the ferromagnetic material. This energy loss is in the form of heat. Therefore, it
is desirable that materials used in electric generators, motors and transformers should have a tall but narrow
hysteresis loop for minimal losses.
19.8 eneRgy LOSS Due tO hySteReSIS LO6

During the process of magnetisation, a loss of energy is always involved in aligning the domains in the
direction of the applied magnetic field. When the direction of an external magnetic field is reversed, the
absorbed energy is not completely recovered and rest energy in sample is lost in the form of heat. This loss
of energy is called hysteresis loss.
19.8.1 Calculation of hysteresis Loss

It can be proved that the energy lost per unit volume of the substance in a complete cycle of magnetisation is
equal to the area of the hysteresis loop. We consider a unit volume of the ferromagnetic substance, which has
N magnetic domains. Let M be the magnetic moment of each magnetic domain which makes an angle q with
the direction of the magnetising field H.
So, the total magnetic moment per unit volume in the direction of magnetising field
= Â M cos q
N
= Intensity of magnetisation (I) (i)
Since there is no magnetisation perpendicular to H, the total magnetic moment perpendicular to the
magnetising field (H),
= Â M cos q = 0 (ii)
N
\ I = Â M cos q
N
or dI = - Â M sin q dq (iii)
N
Here the negative sign shows that I decreases with increasing q.

The work done in rotating the domain from this direction by an angle dq towards the H is
dW = - m0 Â MH sin q dq
N
= m0HdI [By using Eq. (iii)]

The work done per unit volume of substance in the complete cycle is (iv)
Ú m0 H dI
W =
Ú H dI
= m0 
= m0 ¥ (area of I-H loop) (v)
Hence, the work done per unit volume of the substance per cycle of magnetisation is equal to m0 times the
area of I–H curve and this energy is lost in the form of heat.
19.8.1.1 Hysteresis Loss due to B–H curve
The magnetic flux density (B) in substance is due to the magnetising field (H) and the intensity of magnetisation I.
They are related as
B = m0(H + I) (vi)
or dB = m0(dH + dI)
1
or dI = dB - dH (vii)
m0
From Eqs. (v) and (vii) we get
Ú H dB - m0 Ú H dH
W = (viii)
The value of  Ú H dH will be zero, because the curve between H and H is a straight line and will not enclose
any area, i.e.,
Ú H dH = 0
Then Eq. (viii), takes the form,
Ú H dB = area of B–H curve.

W =
19.9 IMPORtance Of hySteReSIS cuRve LO6

The hysteresis curves (B–H loops or I–H loops) of a ferromagnetic material provide very useful information
regarding the magnetic properties of a material. The size and shape of the hysteresis loop of ferromagnetic
materials can be used to obtain information regarding retentivity, coercivity, susceptibility, permeability and
loss in energy per cycle of magnetisation. On the basis of these properties of the magnetic materials, it is
possible to choose a particular material for a particular purpose.
19.10 MagnetIc cIRcuItS LO7

The concept of a magnetic circuit is based on solving some magnetic field problems using circuit approach.
Transformers, motors, generators and relays are magnetic devices which may be considered as magnetic
circuits. By exploiting an analogy between a magnetic circuit and an electric circuit, the analysis of such
circuits is made simple. However, these two types of circuits have some differences. For example, magnetic
flux does not flowÆ
unlike an electric circuit where the current I flows. Further, the permeability m varies with
Æ
the flux density B in a magnetic circuit, whereas the conductivity s is independent of the current density J
in an electric field. This is because ferromagnetic (nonlinear) materials are normally used in most practical
magnetic devices.
19.11 fORceS On MagnetIc MateRIaLS

Consider an electromagnet made of iron of constant relative
permeability, whose coil has N turns and carries a current I, as I
shown in Fig. 19.8. Further, we take that the magnetic field in
the air gap dl is the same as in iron (B1n = B2n).
N Turns
Under the above situation, the attractive force between the two
pieces of iron can be obtained as
dl 1/2 F F 1/2 F
Ê B2 S ˆ
F = - 2Á
Ë 2 m ˜¯
0
Figure 19.8
where S is the area of the coil. The force F is exerted on the
lower piece and not on the current-carrying upper piece giving rise to the field. So the tractive force across
a single gap would be
B2 S
F=-
2 m0
where B is the magnetic field at the surface of the material. The tractive pressure is nothing but the energy
density in the air, given by
B2 1
P= = BH
2m0 2
19.12 MagnetIc MateRIaLS anD theIR aPPLIcatIOnS LO8

Electrical devices like power transformers, motors, generators, electromagnets, etc., use soft magnetic
materials. Electrical steels are used as core materials in them. For retaining magnetic fields of permanent
magnets, hard magnetic materials are used in fabrications. Different types of magnetic materials are used for
different applications.
19.12.1 Low–carbon Steel
Pure iron, although has higher permeability, causes more eddy current losses due to its higher electrical conductivity.
Low–carbon steel has relatively small permeability and higher resistivity. It is the lowest-grade core material.
19.12.2 Iron–Silicon alloys
Adding of about 3–4% silicon to iron produces iron–silicon alloys with improved characteristics. Silicon
increases the electrical resistivity of low-carbon steel and thus reduces the eddy current losses. It also increases
the magnetic permeability and lowers hysteresis losses. It reduces the magnetorestriction and therefore reduces
transformer noise. However, iron–silicon alloys are not useful for communication applications due to their
low magnetic permeability at low fields, because in communication applications much higher permeabilities
are required at low fields.
19.12.2.1 Grain Orientation
By using favourable grain orientation in the material, the hysteresis losses can be decreased and permeability
can be substantially increased. The <100> direction is the easy direction in the case of iron crystals and
spin moments in a virgin crystal are aligned along <100> directions. The <100> direction is parallel to the
rolling direction when steel of iron alloys are manufactured by rolling and annealing. Thus, cold rolled grain
orientated (CRGO) steel carries better magnetic properties in the same direction as that of the direction of
rolling. Consequently, less material is required for cores.
19.12.3 nickel–Iron alloys
If a nickel content of about 25% is present, a pure nickel–iron alloy is practically nonmagnetic. Wide ranges
of magnetic properties are obtained by increasing the nickel content. So, nickel–iron alloys are used for these
applications. Based on the content of nickel, these alloys are divided into three groups: 36% nickel, 50%
nickel and 77% nickel. 36% nickel alloys have high resistivity and low permeability and are used for high–
frequency devices such as speed relays, wideband transformers and inductors. Having moderate permeability
of about 25,000 and high saturation induction, the 50% nickel alloys are used where low loss and small size
are required, such as in small motors, synchores, etc. The 79% nickel alloys have high permeability but lower
saturation induction and are used in recording heads, pulse transformers, sensitive relays, etc.
19.12.4 Mumetal
Multicomponent nickel–iron alloys like permalloy, supermalloy, etc. have the highest permeabilities of the
order of 105. Mumetal having a component of 77% nickel, 16% iron, 5% copper and 2% chromium can be
rolled into thin sheets and is used to shield electronic equipment from stray magnetic fields.
19.12.5 alnico alloys
Alnico alloys containing Al, Ni, Co and Fe and minor constituents of Cu and Ti are used for making
permanent magnets. These are characterised by a high energy product, a high remanent induction and a
moderate coercivity. Besides being mechanically hard and breakable, magnetic properties of alnico alloys
are highly stable against temperature variation, shock, etc. The properties are improved by heat treatments in
alnico 2 or by cooling the alloy in magnetic field.
19.12.6 other alloys
Rare earth magnetic alloys like Sm–Co alloys are superior to alnico alloys in terms of magnetic properties
and they have an energy product up to 2.4 ¥ 105 J / m3 and coercivities of about 3.2 ¥ 106 A / m. They are used
in medical devices such as thin motors in implantable pumps and valves.
Fe–Cr–Co alloys, which are similar to alnico alloys, are used in making permanent magnets for modern
telephone receivers. Similarly, Nd–Fe–B magnetic alloys have a high energy product of the order of
3 ¥ 105 J / m3 and are used mainly in making light and compact electric motors.
19.12.7 Soft ferrites
Due to high electrical properties of dielectrics combined with the magnetic properties of ferromagnetic
materials, ferrites can be used for high frequency applications without eddy current losses. They also have
high electrical resistance (105 to 1015 times the resistance of metallic ferromagnets). The soft ferrites are
used for low signal, memory core, audiovisual and recording head applications. Major applications include
deflection yoke cores, flyback transformers and convergence coils for television receivers. Mn–Zn ferrites
are used for operations of up to 500 kHz whereas Ni–Zn ferrites are effective for the use for high frequency
operation up to 100 MHz.
Mg–Mn ferrites, Mn–Cu ferrites and Li–Ni ferrites are used as memory or logic operation devices in
computers, as switching devices, and in information storage. They are made in the form of tiny rings called
cores, which are assembled into large matrix software containing cores at each junction. Microwave devices
like modulators, couplers, circulators, phase shifters, matching devices are made using microwave ferrites,
mainly manganese ferrite, nickel ferrite, cobalt ferrite, etc.
19.12.8 hard ferrites
Hard ferrites are also used in making permanent magnets. Barium ferrites (trade name Ferroxdure) are being
replaced by strontium ferrites having superior magnetic properties. They find major applications in generator
relays, loudspeakers, telephone ringers, toys, etc. Hard ferrite powders are often mixed with plastic materials
to form flexible magnets for door closers and other holding devices.
19.12.9 Magnetic Storage
Magnetic materials find significant use in the storage of information. Credit cards are popularly used which
also have magnetic strips. To store larger quantities of information at low cost, computers are usually backed
up with magnetic disks.
The recording head consisting of a laminated electromagnet is made of permalloy or soft ferrite having 0.3 m
wide air gap. Here the data written by the electrical signal generates a magnetic field across the gap within
the coil. Finally, the stored information is read using the same head, and an alternating e.m.f is induced in
the coil of the head by moving tap or disk in the read or playback mode. This e.m.f is amplified and fed to a
suitable output device.
S UMMARY
The topics covered in the chapter are summarised below.
✦ Based on their permeability or susceptibility, the magnetic materials can be broadly classified into three
categories, namely diamagnetic, paramagnetic and ferromagnetic materials. The magnetic materials for
which susceptibility cm is negative (permeability mr £ 1) are said to be diamagnetic materials, whereas
the materials with positive susceptibility (permeability mr ≥ 1) are said to be paramagnetic materials.
If the susceptibility is much larger than zero and permeability mr >> 1, then the magnetic materials are
called ferromagnetic materials.
Æ Æ
✦ Paramagnetic and diamagnetic
Æ
materials
Æ
have a linear relationship between B and H . However,
the relationship between B and H is nonlinear for ferromagnetic materials. A material is said to be
nonmagnetic if susceptibility cm = 0 (or mr = 1), it is magnetic otherwise.
✦ The magnetic materials are of two types, namely soft materials and hard materials. Soft magnetic
materials are used in ac applications, since they are easily magnetised and demagnetised. However,
hard magnetic materials are used in producing permanent magnets, since they retain magnetism on a
permanent basis. Due to such properties, these materials are significantly used in information storage
devices.
✦ The magnetic properties of solids originate due to the motion of electrons. The magnetic moment M
neh
of an electron is given by M = , where n = 1, 2, 3…, e is the electronic charge, m is the electron
4p m
mass and h is the Planck’s constant.
✦ When placed in an external magnetic field, the materials which acquire feeble magnetism in the
opposite direction to that of the applied field are called diamagnetic materials. The substances whose
outermost orbits have an even number of electrons show the property of diamagnetism. Bismuth, zinc,
copper, silver, gold, lead, water, etc., are the examples of diamagnetic materials.
✦ When placed in an external magnetic field, the materials which acquire feeble magnetism in the
direction of an applied field are called paramagnetic materials. The source of paramagnetism is the
permanent magnetic moment possessed by the atoms of paramagnetic materials. Aluminium, odium,
platinum, manganese, copper chloride, liquid oxygen, etc., are the examples of paramagnetic materials.
✦ When placed in an external magnetic field, the materials which acquire strong magnetism in the
direction of an applied field are called ferromagnetic materials. This property is found in the substances
which are generally like paramagnetic materials. These are strongly attracted by magnets. Iron, nickel,
cobalt, magnetite (Fe3O4), etc., are the examples of ferromagnetic materials.
✦ Anti-ferromagnetic substances are crystalline materials, in which the dipole moments of the
neighbouring dipoles are equal and opposite in the orientation so that the net magnetisation vanishes.
If they are placed in the magnetic field, they are feebly magnetised in the direction of the field. The
susceptibility of these materials varies with temperature. It increases with increasing temperature and
reaches a maximum at a particular temperature called the Neel temperature (TN). MnO, FeO, CaO,
MnO4, MnS, etc., are the examples of anti-ferromagnetic materials.
✦ If the spins of the atoms are such that there is a net magnetic moment in one direction, the materials are
called ferrimagnetic materials. The examples of ferrimagnetic materials are ferrites which consist of
mainly ferric oxide Fe2O3, combined with one or more oxides of divalent metals.
✦ The classical theory of diamagnetism was developed by the French Physicist Paul Langevin in 1905,
m NZe 2 R 2
according to which the magnetic susceptibility of a diamagnetic material is given by c m = - 0 .
6m
Here, N is the number of atoms per unit volume of the substance, R is the average value of radii for all
the electrons along three axes, m0 is the permeability of free space and Z is the atomic number.
✦ According to Langevin’s theory of paramagnetism, the magnetic susceptibility of a paramagnetic
m I 02
material is given by c m = . Here, I0 is the saturation value of the intensity of magnetisation
3 NkT
I when all the magnetic dipoles get aligned in the direction of an external magnetic field, k is the
m I 02
Boltzmann constant, T is the temperature and is called Curie constant. This equation shows
3 Nk
that the magnetic susceptibility of a paramagnetic material depends on the temperature T and it varies
inversely with T.
✦ Langevin’s theory was not able to explain the complicated dependence of susceptibility on the
temperature, as shown by several paramagnetic substances. In view of this, Langevin’s theory
was further modified by Curie and Weiss. Weiss assumed that in a paramagnetic substance an
internal molecular magnetising field is generated because of a mutual interaction between the
atomic magnetic dipoles. Finally, the magnetic susceptibility of a paramagnetic material is given
C NM 2 ml
by c m = , where q = is called Curie temperature together with M as the permanent
T -q 3k
NM 2 m
magnetic moment of the atoms and l as the molecular field coefficient. C = is the Curie
constant. 3k
✦ In general, a specimen of a ferromagnetic substance contains a number of small regions called domains.
According to the classical theory of ferromagnetism, every domain is magnetically saturated and the
direction of magnetisation in different domains is different. In the absence of an external magnetic
field, all the magnetic domains are randomly oriented and hence their resultant magnetic moment
in any direction will be zero. As per modern theory, it is assumed that in the absence of an external
magnetic field, these domains form closed loops within the substance so that the net magnetic moment
of the whole substance is zero. When a magnetic field is applied to the ferromagnetic substance, the
substance becomes magnetised. When the substance is placed in a weak external magnetic field,
the magnetisation produced is due to the displacement of boundaries of domains and if the external
magnetic field is strong, the magnetisation produced is mainly by the rotation of domains.
✦ Ferromagnetic materials like iron and steel are used for screening or shielding that protects sensitive
electrical devices from disturbances from strong magnetic fields. For perfect screening, it is required
that the shield has infinite permeability (mr = •).
  
✦ Even though
Æ
B =Æ m0 ( H + I ) holds good for all materials including ferromagnetics, the relationship
between B and H depends on previous magnetisation ofÆ a ferromagnetic
Æ Æ
material,
Æ
i.e., its magnetic
history. Instead of having a linear relationship between B and H (i.e., B = mH ), it is only possible to
represent the relationship by a magnetisation curve or B–H curve. Hysteresis is defined as the lagging
of intensity of magnetisation from the magnetising field.
✦ During the process of magnetisation, a loss of energy is always involved in aligning the domains in the
direction of the applied magnetic field. When the direction of an external magnetic field is reversed, the
absorbed energy is not completely recovered and the rest of the energy in the sample is lost in the form
of heat. This loss of energy is called hysteresis loss. The energy lost per unit volume of the substance
in a complete cycle of magnetisation is equal to the area of the hysteresis loop.
✦ The concept of a magnetic circuit is based on solving some magnetic field problems using the circuit
approach. Transformers, motors, generators and relays are magnetic devices which may be considered
as magnetic circuits. By exploiting an analogy between magnetic circuits and electric circuits the
analysis of such circuits is made simple.
✦ Magnetic materials find diverse applications in various modern technologies. Different types of
magnetic materials are used for different applications. Electrical devices like power transformers,
motors, generators, electromagnets, etc. use soft magnetic materials. Electrical steels are used as core
materials in them. For retaining magnetic field of permanent magnets, hard magnetic materials are used
in fabrications.
✦ In view of the application of magnetic materials, different materials were discussed, viz., low-carbon
steel, iron–silicon alloys, nickel–iron alloys, mumetal, alnico alloys, soft ferrites, hard ferrites, etc.
S OLVED E XAMPLES
E xamplE 1 In hydrogen atom, an electron revolves around a nucleus in an orbit of 0.53 Å radius. If the
frequency of revolution of an electron is 6.6 ¥ 1015 Hz, find the magnetic moment of the orbiting electron and
calculate numerical value of Bohr magneton.
Solution Given r = 0.53 ¥ 10–10 m and n = 6.6 ¥ 1015 Hz.
Magnetic Moment M = iA
e e
i= = = en = 1.6 ¥ 10-19 ¥ 6.6 ¥ 1015 A
T 1
n
Area = pr2 = 3.14 ¥ (0.53 ¥ 10–10)2
\ M = iA = 1.6 ¥ 10–19 ¥ 6.6 ¥ 1015 ¥ 3.14 ¥ (0.53 ¥ 10–10)2
= 9.314 ¥ 10–24 Am2
Bohr magneton is the smallest value of the orbital magnetic moment of the electron. For n = 1, Bohr magneton
eh
mB =
4p m
1.6 ¥ 10-19 ¥ 6.6 ¥ 10-34
=
4 ¥ 3.14 ¥ 9.1 ¥ 10-31
= 9.239 ¥ 10-24
24
9.24 10 J /T
E xamplE 2 Determine the magnetisation and flux density in silicon, if its magnetic susceptibility is –4.2 ¥ 10–6
and the magnetic field in it is 1.19 ¥ 105 Am–1. What would be the value of the relative permeability of the
material?
Solution Given c = –4.2 ¥ 10–6 and H = 1.19 ¥ 105 Am–1.
The formulae used are
Magnetisation I = cH
= –4.2 ¥ 10–6 ¥ 1.19 ¥ 105 Am–1
= –0.4998 Am–1
= –0.50 Am–1
Flux density B = m0(H + I)

= 4p ¥ 10–7 ¥ (1.19 ¥ 105 – 0.50)
= 0.1495 T
= 0.150 T
Relative permeability
Ê Iˆ
B = m H = m0 H Á1 + ˜
Ë H¯
m I
or mr = =1+
m0 H
-0.50
=1+
1.19 ¥ 105
= 1 - 0.42 ¥ 10 -5
0.999
E xamplE 3 Find the percentage increase in magnetic induction when the space within a current-carrying
toroid is filled with magnesium. Given c for magnesium as 1.2 ¥ 10–5.
Solution Given c = 1.2 ¥ 10–5.
Magnetic flux density
B = m0 H (i)
When the free space is filled with magnesium, then
B´ = mrm0 H (ii)
and mr = 1 + c (iii)
From Eqs. (ii) and (iii)
B´ = (1 + c)B (iv)
Hence, the percentage increase in magnetic induction
B¢ - B
= ¥ 100 (v)
B
By using Eqs. (i) and (iv), Eq. (v) becomes
(1 + c ) B - B
= ¥ 100 = c ¥ 100
B
= 1.2 ¥ 10–5 ¥ 100
= 1.2 ¥ 10–3%
= 0.0012%
E xamplE 4 Determine the magnetisation and flux density of the diamagnetic material if its magnetic
susceptibility is –0.4 ¥ 10–5 and the magnetic field in it is 104 Am–1.
Solution Given c = –0.4 ¥ 10–5 and H = 104 A m–1.
Magnetisation
I = cH
= –0.4 ¥ 10–5 ¥ 104
= –0.04 Am–1
Magnetic flux density
B = m0(H + I)
= 4p ¥ 10–7 ¥ [104 – 0.04]
= 0.01256 T
E xamplE 5 The magnetic susceptibility of aluminium is 2.3 ¥ 10–5. Find its permeability and relative
permeability.
Solution Given c = 2.3 ¥ 10–5.

Permeability
m = m0mr and
Relative permeability
mr = 1 + c
\ mr = 1 + 2.3 ¥ 10–5
mr = 1.000023
m = m0mr = 4p ¥ 10–7 ¥ 1.000023
= 12.56 ¥ 10–7 N/A2
E xamplE 6 The magnetic susceptibility of a medium is 940 ¥ 10–4. Calculate its absolute and relative
permeability.
Solution Given c = 0.094.
Absolute permeability
m = m0mr
Relative permeability mr = 1 + c
so, mr = 1 + 0.094 = 1.094
m = m0mr = 4p ¥ 10–7 ¥ 1.094
= 13.74 ¥ 10–7 N/A2
E xamplE 7 The maximum value of the permeability of a material is 0.126 N / A2. What is the relative
permeability and magnetic susceptibility?
Solution Given m = 0.126N / A2.
m
Relative permeability m r = and susceptibility is c = mr – 1
m0
m 0.126
\ mr = =
m0 4p ¥ 10-7
= 105
c = m r - 1 = 105 - 1
99999
E xamplE 8 Calculate the diamagnetic susceptibility of He assuming that the two electrons are contributing
to its diamagnetism. Consider the mean radius of the atom as 0.6 Å and N = 28 ¥ 1026 per m3.
Solution Given N = 28 ¥ 1026 per m3 and R = 0.6 ¥ 10–10 m.
Susceptibility of diamagnetic material
- m0 Ze 2 NR 2
c dia =
6m
(4p ¥ 10 -7 ) ¥ 2 ¥ (1.6 ¥ 10 -19 )2 ¥ 28 ¥ 1026 ¥ (0.6 ¥ 10-10 )2
=-
6 ¥ 9.1 ¥ 10 -31
8
11.872 10
E xamplE 9 A magnetising field of 1000 A / m produces a magnetic flux of 2 ¥ 10–5 Weber in a bar of iron of
2
0.2cm cross-section. Calculate permeability and susceptibility of the bar.
Solution Given H = 103 A / m, f = 2 ¥ 10–5 Wb and A = 0.2 ¥ 10–4 m2.
f
Magnetic flux density B =
A
B
Permeability m= –
H
m
Susceptibility c = mr - 1 = -1
m0
f 2 ¥ 10-5
\ B= = = 1.0 Wb/m 2
A 2 ¥ 10-5
B 1
m= = = 10-3 N/A 2
H 103
m 10-3
and c= -1= - 1 795.18
m0 4p ¥ 10-7
E xamplE 10 An iron rod of 1.0 m length and cross-section 4 sq cm is in the form of a closed ring. If the
permeability of iron is 50 ¥ 10–4 Hm–1. Show that the number of ampere turns required to produce a magnetic
flux of 4 ¥ 10–4 Wb through the closed ring is 200.
Solution Given L = 1.0 m, A = 4 ¥ 10–4 m2, m = 50 ¥ 10–4 H/m and f = 4 ¥ 10–4 Wb.
f
Magnetic flux density B =
A
4 ¥ 10-4
= = 1.0 Wb/m 2
4 ¥ 10-4
Also B = mNI
B
\ Ampere turn NI =
m
1.0
= 200 A/m
50 ¥ 10-4
E xamplE 11 The mean length of an iron ring having 200 turns of wire upon it is 0.5 m and its cross-section
is 4 ¥ 10–4 m2. What current through the winding should be sent to produce a flux of 4 ¥ 10–4 Wb in the ring?
Permeability of iron is 65 ¥ 10–4 Wb / Am.
Solution Given m = 6.5 ¥ 10–4 Wb / Am, f = 4 ¥ 10–4 Wb and A = 4 ¥ 10–4 m2.
f 4 ¥ 10-4
The formula used is B = = = 1.0 Wb/m 2
A 4 ¥ 10-4
Also B = m NI
B
or I=
mN
where N is the number of turns per metre, i.e.,
200
N= = 400 turns/m
0.5
B 1.0
Then, current I = =
m N 6.5 ¥ 10-4 ¥ 400
I 3.85 A
E xamplE 12 Assuming the susceptibility of a diamagnetic material as –5.6 ¥ 10–6 and its structure as a body-
centred cubic with a lattice constant 2.55 Å.
Calculate the radius of its atom, if only one electron per atom is contributing to diamagnetism.
Solution Given c = –5.6 ¥ 10–6 and a = 2.55 Å = 2.55 ¥ 10–10m.
m0 Ze 2 NR 2
c=-
6m
1/ 2
È - c 6m ˘
or R=Í 2 ˙
Î m0 Ze N ˚
where N is the number of atoms per unit volume
1 1 1
i.e., N =2 =2¥ 3 =2¥
V a (2.55 ¥ 10-10 )3
where the factor 2 arises because the body-centred cubic has two electrons per unit cell.
or N = 1.206 ¥ 1029 per m3
1
È 5.6 ¥ 10-6 ¥ 6 ¥ 9.1 ¥ 10-31 ˘2
\ | R| = Í -7 -19 2 29 ˙
Î 4p ¥ 10 ¥ 1 ¥ (1.6 ¥ 10 ) ¥ 1.206 ¥ 10 ˚
0.89 Å
E xamplE 13 A paramagnetic substance contains 6.5 ¥ 1025 atoms per m3 and the magnetic moment of each
atom is one Bohr magneton. Find the susceptibility at room temperature.
Solution Given N = 6.5 ¥ 1025 atoms / m3.
T corresponding to room temperature = 27 + 273 = 300 K
m0 NM 2
Susceptibility c= (i)
3kT
The magnetic moment of each atom
Ê eh ˆ
M = nÁ
Ë 4p m ˜¯
Ê 1.6 ¥ 10-19 ¥ 6.6 ¥ 10-34 ˆ
= 1Á ˜
Ë 4 ¥ 3.14 ¥ 9.1 ¥ 10-31 ¯
= 9.24 ¥ 10-24 Am 2
From Eq. (i), we get
(4p ¥ 10-7 ) ¥ 6.5 ¥ 1025 ¥ (9.24 ¥ 10-24 )2
c=
3 ¥ 1.38 ¥ 10-23 ¥ 300
= 5.612099 ¥ 10-7
7
5.612 10
E xamplE 14 The molecular weight and density of a paramagnetic substance are 168.5 and 4370 kg / m3,
respectively, at room temperature. Considering the contribution to paramagnetism as two Bohr magnetons
per molecule, calculate its susceptibility and magnetisation produced in it in a field of 2 ¥ 105 Am–1.
Solution Given molecular weight M0 = 168.5, molecular density f = 4370 kg / m3, room temperature (T) = 27°C
= 27 + 273 = 300 K, H = 2 ¥ 105 A / m, Bohr magnetons (mB) = 9.24 ¥ 10–24 A m2.
fN A 4370 ¥ 6.02 ¥ 1026

Number of atoms per unit volume N = =
M0 168.5
= 1.56 ¥ 1028 per m3
Susceptibility c = m0Nm2B / 3kT

(4p ¥ 10-7 ) ¥ (1.56 ¥ 1028 ) ¥ (2 ¥ 9.24 ¥ 10-24 )2
=
3 ¥ (1.38 ¥ 10-23 ) ¥ 300
In the above expression factor of 2 comes before mB, as the contribution is two Bohr magnetons per molecule
\ cm = 5.388 ¥ 10–4
Now magnetisation I = cH
= 5.388 ¥ 10–4 ¥ 2 ¥ 105
= 107.76 A / m
E xamplE 15 The hysteresis loop of a transformer has an area of 2500 ergs / cm3. Calculate the loss of energy
per hour at 50 Hz frequency. The density of iron is 7.5 g / cm3 and weight is 10 kg.
Solution Given m = 10 kg = 104 g and d = 7.5 g / cm3.
Area of hysteresis loop = 2500 ergs / cm3
The loss of energy per unit volume per hour
= 50 ¥ 60 ¥ 60 ¥ 2500
= 4.5 ¥ 108 ergs / cm3
m 104
Volume of iron (V) = = = 1 ◊ 33 ¥ 103 cm3
d 7◊5
Hence, the total loss of energy per hour
= 4.5 ¥ 108 ¥ 1.33 ¥ 103
= 6.0 × 1011 ergs
E xamplE 16 A bar magnet has a coercivity of 5 ¥ 103 A / m. It is desired to demagnetise it by inserting it inside
a 10 cm long solenoid having 50 turns. What current should be sent through the solenoid?
Solution Here, coercivity H = 5 ¥ 103 A / m, l = 10 cm and total twins = 50,
Turns per meter N = 50 ¥ 10 = 500 turns / m
Now, H = Ni
5 ¥ 103 = 500 ¥ i
or i = 10 A
ExamplE 17 An iron rod of 50 cm length and 4 sq cm cross-section area is in the form of a circular ring. If the
permeability of iron is 65 ¥ 10–4 H/m, compute the number of turns required to produce a flux of 4 ¥ 10–5 Weber.
Solution Here, A = 4 cm2 = 4 ¥ 10–4 m2, l = 50 cm = 0.5 m,

m = 65 ¥ 10–4 H / m and f = 4 ¥ 10–5 Wb.
N
As H = i or Ni = Hl
l
f 4 ¥ 10-5
B= = = 0.1 Wb/m 2
A 4 ¥ 10-4
B B 0.1
Q m= H= = = 15.38 A-turns/m
H m 65 ¥ 10-4
Number of turns = N = Hl = 15.38 ¥ 0.5 = 7.69
E xamplE 18 A magnetising field of 600 Am–1 produces a magnetic flux of 2.4 ¥ 10–5 Weber in an iron bar of
2
0.2 cm cross-sectional area. Compute the permeability and susceptibility of the bar.
Solution Given f = 2.4 ¥ 10–5 Wb and A = 0.2 cm2 = 0.2 ¥ 10–4m2. The magnetic flux density is given by
f 2.4 ¥ 10-5
B= = 1.2 wb /m 2
A 0.2 ¥ 10-4
The permeability is given by
B
m=
H
1.2
or m= 0.002 N/A 2
600
The susceptibility is given by
m
c= -1
m0
0.002
= - 1 = 1592 - 1
4 ¥ 3.14 ¥ 10-7
1591
E xamplE 19 The magnetic susceptibility of medium is 950 ¥ 10–11. Compute the permeability and relative
permeability.
Solution Magnetic susceptibility c = 950 ¥ 10–11.
As m = m0(1 + c)
and permeability of free space m0 = 4p ¥ 10–7 H / m
\ m = 4p ¥ 10–7 ¥ (1 + 950 ¥ 10–11)
Hence, m is slightly greater than m0
11
Now relative permeability r 1 950 10
0
E xamplE 20 Find the energy loss per hour in an iron core of a transformer, if the area of the B–H loop is
250 J / m3 and the frequency of the alternating current is 50 Hz. The density of iron is 7.5 ¥ 103 kg / m3 and the
mass of the core is 100 kg.
Solution Area of B–H loop = 250 J / m3 and frequency f = 50 Hz, density r = 7.5 ¥ 103 kg / m3and mass m = 100 kg.
m 100
Volume of core, V = = = 13.3 ¥ 10-3 m3
r 7.5 ¥ 103
Number of ac cycles in 1 hour = (60 ¥ 60) ¥ 50

= 180000
Hysteresis loss per cycle per unit volume = area of B–H loop = 250 J / m3
\ Hysteresis loss in the entire core in 1 hour
= (250) ¥ 13.3 ¥ 10–3 ¥ 180000 = 59.85 × 104J
E xamplE 21 In B–H loop, the maximum value of Bmax is 1.375 Weber / m2 and the area of the loop is
0.513 cm . If the value of 1 cm on the x-axis is 10 A / cm and the value of 1 cm on the y-axis is 1 Weber / m2,
2
calculate the hysteresis power loss when an alternating magnetic flux density of 1.375 Weber / m2 intensity
and 50-Hz frequency is applied on 10–3m3 volume of the specimen.
Solution 1 cm on the x-axis = 10 A / cm = 10 ¥ 100 A / m
1 cm on the y-axis = 1 Wb / m2
\ Area of the B–H loop = 0.513 cm2 = 0.513 ¥ (10 ¥ 100) ¥ 1 = 513 J / m3
Hysteresis loss per cycle per m3 = 513 J
But volume of specimen = 10–3m3
and frequency = 50 Hz ( = number of cycles per second)
\ Hysteresis loss per second (or hysteresis power loss)
= 513 ¥ 10–3 ¥ 50
= 25.65 W
Q.1 Which of the following statements is/are true about magnetic susceptibility c?
(a) c may be positive or negative
(b) c for a paramagnetic material has values close to 1
(c) At a given temperature, the value of cm increases with increasing magnetic field
(d) For paramagnetic substances, c is inversely proportional to the absolute temperature of the same
Q.2 The relative permeability of a medium is the permeability relative to that of
(a) vacuum (b) iron (c) water (d) none of these
Q.3 The dipole moment of current loop does not depend upon the
(a) current in the loop (b) shape of the loop
(c) area of the loop (d) number of turns in the loop
Q.4 The dimensions of magnetic susceptibility are
(a) Wb / m (b) amp / m (c) Wb / m2 (d) dimension less
Q.5 The magnetic susceptibility of a diamagnetic substance is
(a) positive (b) negative (c) zero (d) none of these
Q.6 The magnetic susceptibility is negative for a substance; which is:
(a) diamagnetic (b) ferromagnetic (c) paramagnetic (d) none of these
Q.7 Which of the following materials are feebly attracted by external magnetic fields?
(a) ferromagnetic material (b) ferrimagnetic material
(c) parramagnetic material (d) none of these
Q.8 The magnetic susceptibility of a diamagnetic substance

(a) does not depend on temperature (b) increases with increase in temperature
(c) decrease with decrease in temperature (d) none of these
Q.9 Which of the following substances are diamagnetic material?
(a) alluminium and platinum (b) bismuth and lead
(c) copper and gold (d) both (b) and (c)
Æ
Q.10 A current loop of magnetic moment P m when placed in a non-uniform magnetic field experiences a
force
Æ Æ Æ Æ Æ Æ Æ Æ Æ Æ Æ
(a) —(P m ◊ B ) (b) (— ◊ B ) (c) — ◊ P mB (d) (P m ◊ —)B
Q.11 The relation between magnetising field H and intensity of magnetisation is
(a) I µ H–1 (b) I µ H (c) I µ H2 (d) I µ H1/2
Q.12 In classical theory of diamagnetism, the term eB/2 m is known as
(a) Bohr magneton (b) Debye angular frequency
(c) Larmour angular frequency (d) none of these
Q.13 Numerical value of Bohr magneton is
(a) 9.24 ¥ 10–24 J/T (b) 6.6 ¥ 10–24 J/T (c) 6.6 ¥ 10–14 J/T (d) none of these
Q.14 The value of susceptibility of a diamagnetic substance is about
(a) 105 (b) 10–6 (c) –10–7 (d) 107
Q.15 The magnetic moment of an atom is due to
(a) orbital motion of electrons only (b) spin motion of electron
(c) both orbital and spin motion (d) none of these
Q.16 The value of susceptibility of a ferromagnetic substance is of the order of
(a) 105 (b) 106 (c) 10–7 (d) –10–6
Q.17 Which of the following substance/substances can have positive permeability and negative susceptibility?
(a) diamagnetic (b) ferromagnetic (c) paramagnetic (d) none of these
Q.18 Feritites show
(a) paramagnetic behaviour (b) ferrimagnetic behaviour
(c) ferromagnetic behaviour (d) none of these
Q.19 On application of magnetic field on a material, if it does not show any effect, then the material is a
(a) non-magnetic material (b) diamagnetic material
(c) anti-ferromagnetic material (d) none of these
Q.20 The boundary wall between domains is known as
(a) potential wall (b) Bloch wall (c) magnetic wall (d) none of these
Q.21 Which of the following statement is correct?
(a) Magnetic dipole moment per unit volume of material is known as magnetisation of substance
(b) In a ferromagnetic material, susceptibility is very small and positive
(c) The most suitable material for making permanent magnet is steel
(d) all of these
Q.22 Which of the following statement is true?
(a) Hysteresis is lagging of an effect behind the cause of effect
(b) In soft magnetic material, the hysterisis loss is less
(c) The area of hysteresis loop of a ferromagnetic gives the energy that is consumed during each cycle
(d) all of these
Q.1 What do you mean by magnetisation, permeability and susceptibility of a magnetic substance?
Q.2 What is the magnetic dipole moment associated with a loop-carrying current?
Q.3 Define atomic magnetic moment and discuss orbital diamagnetism.
Q.4 What is diamagnetism?
Q.5 Discuss ferromagnetism.
Q.6 What do you mean by ferromagnetic domain?
Q.7 How do you account for the magnetic properties of materials? Explain.
Q.8 Is it meaningful to say that an atom is ferromagnetic?
Q.9 Why is ferromagnetism found in solids only and not in fluids?
Q.10 What are the characteristics of diamagnetic, paramagnetic and ferromagnetic substances?
Q.11 What is Curie Point or Curie temperature? The magnetic behaviour of magnetic substances decreases
with increasing temperature. Comment.
Q.12 Explain the temperature dependence of the behaviour of paramagnetic, diamagnetic and ferromagnetic
substances.
Q.13 What does the area of a B–H loop represent?
P RACTICE P ROBLEMS
general Questions
Q.1 Explain magnetic flux density (B), intensity of magnetisation (M), magnetic flux density (H). How are
they related to each other?
Q.2 Define magnetic susceptibility (c) and relative magnetic permeability (mr) and establish a relation
m = m0 (1 + c).
Q.3 Distinguish between dia, para and ferromagnetic materials. Derive an expression for magnetic
susceptibility of a paramagnetic substance.
Q.4 Differentiate paramagnetic, diamagnetic and ferromagnetic substances by illustrating simple experiments.
Q.5 Derive an expression for diamagnetic susceptibility on the basis of Langevin’s theory and show that it
is independent of temperature.
Q.6 Discuss diamagnetic, paramagnetic, ferromagnetic, anti-ferromagnetic and ferrimagnetic substances
citing one example of each.
Q.7 Prove that the change is the same whether the electron is orbiting around the nucleus in clockwise
direction or anti-clockwise direction. Hence, discuss the diamagnetic behaviour of the substance
according to Langevin’s theory of diamagnetism.
Q.8 Based on Langevin’s theory of diamagnetism, show that the diamagnetic susceptibility is negative and
independent of temperature and field strength.
Q.9 Give Langevin’s electronic theory of paramagnetism and hence prove that susceptibility (c) of
paramagnetic substance is inversely proportional to absolute temperature.
Q.10 Why are some substances diamagnetic while others paramagnetic? Explain.
Q.11 How do you classify a material as dia, para or ferromagnetic? Discuss the classical theory of paramag-
netism.
Q.12 Explain the origin of atomic dipole moments and derive Langevin’s equation for paramagnetic
susceptibility.
Q.13 What are the physical basis of diamagnetism and paramagnetism of materials? Describe the Weiss’s
molecular theory of ferromagnetism and derive the Curie–Weiss Law.
Q.14 What are the distinguishing features of ferromagnetism? Give the theory of magnetic domains in
ferromagnetic materials.
Q.15 What is ferromagnetism? Explain ferromagnetism on the basis of domain theory. Why does a piece of
iron ordinarily not behave as a magnet?
Q.16 What do you understand by hysteresis remanence (retentivity) and coercivity? How do you determine
the value of remanence and coercivity from a hysteresis loop?
Q.17 Show that the loss of energy due to hysteresis per unit volume of the material per cycle of magnetisation
is given by (i) m0 ¥ area of I–H loop and (ii) area of B–H loop.
Q.18 What type of material should be used for making
(a) permanent magnets, and
(b) electromagnets?
Q.19 Explain the use of a hysteresis curve. What type of magnetic material is suitable for transformer cores,
telephone diaphragm and chokes?
Superconductivity 20
Learning Objectives
LO 1 Gain knowledge/Learn about electrical LO 5 Discuss penetration depth of
resistivity of solids and phonons supercurrent and magnetic flux in
LO 2 Understand the properties and superconductors
classification of superconductors LO 6 Explain formation of Cooper pairs and its
LO 3 Learn about effect of magnetic field and relation to Bose–Einstein condensation
Isotope effect on superconductivity LO 7 Understand basis of BCS theory and
LO 4 Know how London equations explain coherence length
zero resistance and ideal diamagnetism LO 8 Analyse high temperature conductivity
of superconductors and applications of conductivity
Introduction
The phenomenon of superconductivity was first discovered by Kammerlingh Onnes in 1911. He found that
electrical resistivity of some metals, alloys and compounds drops suddenly to zero when they are cooled
below a certain temperature. This phenomenon is
known as superconductivity and the materials that
exhibit this behaviour are called as superconductors.
Electrical Resistivity
However, all the materials cannot superconduct

even at 0 K. The temperature at which a normal r
cto
material turns into a superconducting state is called n du
o
rc
critical temperature Tc. Each superconducting u pe
S l
material has its own critical temperature. eta
a lM
Kammerlingh Onnes discovered that the electrical rm
No
resistance of highly purified mercury dropped
abruptly to zero at 4.15 K, as shown in Fig. 20.1.
Generally good conductors like Au, Ag, Cu, Li,
O TC Temperature (K)
Na, K, etc. do not show superconductivity even at
Figure 20.1
absolute zero.
The interaction of electrons with one another and with the lattice ions were averaged out by the free
electron theory (model) approximation. This could be responsible for resistance to the flow of electrons
under normal conditions. This independent particle model was unable to explain superconductivity. The
clear understanding of the phenomenon of superconductivity requires the consideration of collective
behavior of electrons and ions. This is called many body effects in solids.
20.1 ELEctricaL rEsistivity Of sOLids and PhOnOns LO1

There are several factors that contribute to the electrical resistivity of a solid. For example, the deviations
from a perfect lattice, which may be due to impurities or structural defects in crystal, can scatter the electrons.
Moreover, the vibrations of lattice ions take place in normal modes. These vibrations constitute acoustic
waves which travel through the solid. These waves are called phonons, which carry momentum. It is obvious
that the number of phonons will increase if the temperature is raised. In the presence of phonons, now there
is an interaction between the electrons and phonons. This interaction scatters conduction electrons and hence
causes more resistance. Therefore, it is clear that the electrical resistance of a solid will decrease if we cool
the solid.
20.2 PrOPErtiEs Of suPErcOnductOrs LO2

Electrical, magnetic and thermal properties are the main properties of the superconductors. These are
described below.
20.2.1 Electrical Property
A superconductor is characterised by zero electrical resistivity. Once the current is started to flow, it will
continue for years without any detectable decay (ideally) even if the applied voltage is removed.
20.2.2 Magnetic Property: Meissner Effect

An important property of the superconducting phase is the repulsion of magnetic flux lines from the bulk of
superconductor. It is called Meissner effect. When a specimen is placed in a magnetic field, the magnetic flux
lines pass through it (Fig. 20.2a). Now, if the temperature is decreased below the transition temperature (Tc),
it expels all the magnetic flux lines from inside of the specimen (Fig.20.2b). Hence, we get
B = m0(H + M) = 0 or M = –H
where M is the intensity of magnetisation due to applied magnetic field H. By the definition of magnetic
susceptibility
M
cm = = -1
H
Since diamagnetic materials have negative magnetic susceptibility, the specimen becomes an ‘ideal
diamagnetic’ in superconducting state.
If a specimen of superconductor is placed in a strong magnetic field, the specimen loses its property of
superconductivity and becomes normal material as shown in Fig. 20.2.
(a) (b) (c)

Figure 20.2
20.2.3 Thermal Properties
Thermal properties include the entropy, specific heat
and thermal conductivity, which are discussed below.
20.2.3.1 Entropy
We know that the entropy is a measure of the disorder Normal State
Entropy
of a system. In all superconductors, the entropy

decreases significantly on cooling below the critical
temperature Tc. Therefore, the observed decrease in
entropy between the normal state and superconducting Superconducting
state shows that the superconducting state is more State
ordered than the normal state. For aluminium, the
change in entropy was observed to be small of the
order of 10−14 k per atom, where k is the Boltzmann TC
constant. The variation of entropy of aluminium in the
O T(K)
normal and superconducting states with temperature is
shown in Fig. 20.3. Figure 20.3
20.2.3.2 Specific Heat
The specific heat of normal metal is found to vary with
temperature. The variation follows the following trend
Superconducting
Cn(T) = gT + bT 3 State
Normal State
In this relation, the first term is the specific heat
[Cv]el
of electrons in metal whereas the second term is

due to the contribution of lattice vibrations at low
temperature. On the other side, the specific heat of
the superconductor shows a jump at Tc (Fig. 20.4). As
we know that the superconductivity affects electrons
TC
mainly, so we may assume that the lattice vibration
part remains unaffected and has the same value bT 3 in O T(K)
the normal and superconducting states. Figure 20.4
20.2.3.3 Thermal Conductivity
The thermal conductivity of superconductors undergoes a continuous change between the two phases. It is
usually lower in the superconducting phase which shows that the electronic contribution goes down. This
suggests that the superconducting electrons possibly play no role in the transfer of heat.
20.3 cLassificatiOn Of suPErcOnductOrs LO2

On the basis of magnetising behaviour, superconductors can be classified as type-I (or soft) and type-II
(or hard) superconductors.
20.3.1 Type-I (Soft) Superconductor

This type of superconductor obeys complete Meissner effect. It expels all the magnetic field abruptly from
the interior and becomes an ‘ideal diamagnetic material’. Magnetisation produced in the superconductor
remains in the direction opposite to the applied external magnetic field, as shown in Fig. 20.5a. At the critical
magnetising field, the magnetisation decreases abruptly and the material becomes normal. For all the values
of external magnetic field above the critical value, the magnetic flux lines penetrate completely inside the
material.
20.3.2 Type-II (Hard) Superconductor

This type of superconductor loses magnetisation gradually rather than abruptly. From the Fig. 20.5b it is clear
that at the magnetic field HC1, the flux starts penetrating into the material until the upper critical field HC2
is reached. Between the two critical magnetic fields HC1 and HC2, the material is said to be in a mixed state.
Above the magnetic field HC2, the material becomes normal conductor. Inspite of the fact that the magnetic
flux lines penetrate inside the material in the mixed state, the electrical resistivity continues to be zero upto
the magnetic field HC2.
Magnetisation (–M)
Magnetisation (–M)
HC HC1 HC2
Applied Magnetic Field Applied Magnetic Field
(a) (b)
Figure 20.5
20.4 EffEcT of MagnETIc fIEld LO3

In order to achieve the superconducting state in the metal, it is required that some combination of temperature
and field strength should be less than a critical value. The superconductivity of a specimen can be destroyed
by applying a critical magnetic field. So a superconductor becomes a normal conductor when it is placed
in an intense magnetic field. Therefore, it retains its resistance. The critical magnetic field is a function of
temperature also and it is defined as
È Ê T ˆ2˘
H c = H 0 Í1 - Á ˜ ˙
ÎÍ Ë Tc ¯ ˚˙
Here H0 is the critical magnetic field at 0 K. We can see that at T = 0 K, Hc = H0, and at T = Tc, Hc = 0. This
relation between Hc and T shows that the critical magnetic field varies parabolically with the temperature.
This curve demarcates the two states, i.e., it defines the boundary below which superconductivity is present
and outside of which it behaves as a normal conductor.
20.5 isOtOPE EffEct LO3

It has been observed that the critical temperature Tc of superconductors varies with isotopic mass. Higher
Tc is found in samples with lighter nuclei. In mercury, Tc varies form 4.85 K to 4.146 K as the average
atomic mass M varies from 199.5 a.m.u. to 203.4 a.m.u. It is also found that the transition temperature
changes smoothly when we mix different isotopes of the same element. The dependence of Tc on the
atomic mass reveals that lattice vibrations and hence electron phonon interaction is deeply involved in the
superconductivity. Based on experimental results it is found that
Tc µ M−a
or TcMa = Constant
Here, M is the atomic mass, Tc is the critical temperature and a = 0.49 ± 0.01. In view of this value of a it
was thought that a = 0.5 is valid for most of the materials. With this we get
TcM1/2 = Constant
or Tc1 M11/ 2 = Tc2 M 21/ 2
20.6 LOndOn EquatiOns LO4

London equations were developed by brothers Fritz London and Heinz London in 1935 in order to explain
the zero resistance and ideal diamagnetism of superconductors. These equations describe superconducting
phenomenon in a simplified and meaningful way by providing relationships between superconducting
current and electromagnetic fields in and around a superconductor. In particular, the London equations
beautifully explain the Meissner effect, wherein a material expels all the internal magnetic fields as it crosses
the superconducting threshold.
London brothers followed a certain intuitive logic in the formulation of their theory. Since the electrons in a
superconductor flow with no resistance, a linear relationship between current and electric field (Ohm's law)
may not be possible in a superconductor. In view of this, London imagined the electrons as if these were free
electrons under the influence of a uniform external electric field. According to the Lorentz force law these
electrons should encounter a uniform force and thus accelerate uniformly. Hence, the equation of motion for
the superconducting electrons in the presence of applied electric field E is given by

dv  
m = F = - eE (i)
dt

The current density J can be expressed as
 
J = –nev (ii)
where n is the number of electrons per unit volume i.e., n is the number density of superconducting carriers.
By differentiating Eq. (ii) w.r.t. time, we get
  
dJ dv È eE ˘
= - ne = - ne Í- ˙
dt dt Î m˚

dJ ne 2 
= E (iii)
dt m
Equation (iii) is known as first London equation. According to London’s theory, it was assumed that two
types of the electrons, i.e., normal and superconducting electrons are present in the superconductors. The
normal electrons don’t respond to the electric field and only the superconducting electrons respond to the
electric field. Now Maxwell’s equation can be written as

 dB
curl E = -
dt

  dH
or — ¥ E = - m0 (iv)
  dt
as B = m0H.
Taking curl of Eq. (iii), we get

dJ ne 2 
curl = curl E
dt m

 dJ ne 2  
or —¥ = (— ¥ E ) (v)
dt m
By using Eqs. (iv) and (v), we get

 
 dJ m0 ne 2 dH
—¥ =- (vi)
dt m dt
By integrating Eq. (vi) w.r.t. time, we get
 m ne 2  
curl J = - 0 [H - H0] (vii)
m
where H0 is a constant of integration. As we know that Meissner effect exhibits complete absence of magnetic
field inside the superconductor. Therefore, H0 must be zero. Then
 m ne 2 
curl J = - 0 H
m
 ne 2 
or curl J = - B (viii)
m
Equation (viii) is known as second London equation which explains Meissner effect as well.
20.7 PEnEtratiOn dEPth LO5

The second London equation (Eq. (viii)) explains the Meissner effect. This equation also predicts the
penetration of supercurrent and magnetic flux in a superconductor, which will be shown below.
We take the curl of Eq. (viii)
   m ne 2    
— ¥ (— ¥ J ) = - 0 —¥H ( B = m0 H )
m
    m ne 2  
—( — ◊ J ) - — 2 J = - 0 —¥H
m
 
Since — ◊ J = 0, we obtain
 m ne 2  
—2 J = 0 —¥H (ix)
m
The Maxwell’s equation for direct current is
   
— ¥ H = J ( x)  ∂E / ∂t = 0 (x)
From Eqs. (ix) and (x), we find
 m ne 2 
—2 J = 0 J
m

2
 J
— J= 2 (xi)
lL
m
where l L2 = (xii)
m0 e2 n
The parameter lL has the dimensions of length and is called London penetration depth. So in terms of the
penetration depth, Eq. (xi) can be written in one dimension as
 
d2J J
=
dx 2 l L2
Solution of this equation is as follows
J = A1e x / l L + A2 e - x / l L
where A1 and A2 are constants. x represents the distance into the metal (superconductor) from the surface.
The first term in the above solution gets increased when x is increased. However, this is contrary to the fact.
Therefore, we neglect coefficient A1 and write the above solution as
J = A2 e - x / l L
Now at x = 0, J has some finite value, which we consider as J0, i.e., J = J0. With this J = J 0 e - x / l L .
or J = J 0 e- x / lL
J 1
When x = lL, the above solutions reads = .
Penetration Depth × 107 (in m)

J0 e
3
Thus, the penetration depth is the depth where the
current drops to 1/e times of its value at the surface.
It is a fundamental length that characterises a 2
lL
superconductor. Therefore, the London equations
provide a characteristic length scale lL, over
which external magnetic fields are exponentially 1
suppressed. In order to understand the physical TC
meaning of the London penetration depth, we
consider a superconductor within free space where a 0 1 2 3 4
constant magnetic field outside the superconductor T(K)
is pointed parallel to the superconducting boundary
Figure 20.6
plane in the z-direction. Then for the x-direction,
which is perpendicular to the boundary, the solution inside the superconductor may be shown to be
Bz(x) = B0e–x/lL. This explains the exponential suppression of the external magnetic field in the superconductor.
The penetration depth depends on temperature and it gets significantly increased as T approaches Tc
(Fig. 20.6). The variation of penetration depth with temperature is according to the following relation
-1/ 2
l (T ) È Ê T ˆ ˘
4
= Í1 - Á ˜ ˙
l (0) ÍÎ Ë Tc ¯ ˚˙
20.8 cOOPEr Pairs LO6

As mentioned earlier, the superconductors offer zero resistance. So, when an electron in a solid passes by
adjacent ions in the lattice, it imparts a momentum to these ions due to Coulomb interaction. It means the ion
starts vibrating and hence phonon gets excited. Now consider a second electron which is subsequently passing
through the moving region of increased positive charge density. It will experience an attractive Coulomb
interaction and can absorb all momentum of vibrating region. Under this situation, the vibrations of ions
are stopped and hence the phonon is absorbed by this second electron. So in this interaction, the momentum
which was imparted by the first electron is taken up by the second electron and hence these electrons undergo
an interaction. This interaction would be an attractive interaction because exchange of momentum takes
place via Coulomb attraction interaction; of course through phonon. According to BCS theory, under certain
conditions this attractive interaction overcomes the force of repulsion between the two electrons. Therefore,
the electrons are loosely bound together. This pair of electrons is called Cooper pair. According to BCS
theory, the binding energy of Cooper pair at absolute zero is ~ 3kTc. This binding energy gets reduced as the
temperature is raised. The binding energy becomes zero when temperature equals the critical temperature Tc.
Hence, a Cooper pair is not bound at temperature ~ Tc.
In view of the above discussion, it is clear that the Cooper pairs are formed due to the electron lattice
interactions. The two electrons of a cooper pair have equal and opposite momenta. They also have opposite
spins, i.e., one electron is spin up and another is spin down. Thus, the bound Cooper pair is a spin zero object
and is a boson. The bound Cooper pairs overlap each other because the Cooper pair wavefunction is very
large (few hundred nanometers in diameter). As we know that when identical bosons overlap each other, then
a large number of bosons condense into the same quantum state (Bose Einstein condensation). The motions
of the Cooper pairs are strongly correlated because all of the Cooper pairs are in the same quantum state.
20.9 BOsE EinstEin cOndEnsatiOn LO6

Based on the fact that bosons are governed by Bose Einstein statistics and they are not constrained by the Pauli
exclusion principle, Einstein in 1924 had pointed out that bosons could condense in unlimited numbers into
a single ground state. Then very late anomalous behaviour of liquid helium was noticed at low temperatures.
Actually a remarkable discontinuity in heat capacity of helium was observed when it was cooled to critical
temperature of 2.17 K. Under this condition, the liquid density drops and a fraction of the liquid attains
a zero viscosity, i.e., it becomes a superfluid. This superfluidity takes place due to the fraction of helium
atoms which condenses to the lowest possible energy. This is called Bose Einstein Condensation, which
can be achieved when the participating particles are identical. This condition of indistinguishability requires
that the deBroglie wavelengths of the particles overlap significantly. This is possible only at extremely low
temperature as the deBroglie wavelengths become longer. This also requires a high density of particles to
narrow the gap between the particles. Cornell and Wieman together with Ketterle received the 2001 Nobel
Prize for achieving Bose Einstein Condensation in dilute gases of alkali atoms.
20.10 BcS THEory: QualITaTIvE ExPlanaTIon LO7

The basis of a quantum theory of superconductivity was led by the classic 1957 paper of Bardeen, Cooper
and Schrieffer. This is now called the BCS theory. It was widely applicable, for example, from He3 atoms
in their condensed phase to type I and type II metallic superconductors. This theory involves the electron
interaction through phonons as mediators. The formulation of BCS theory is based on two experimental
facts viz. the isotope effect and the variation of specific heat of superconductors. The isotope effect indicated
that the lattice vibrations play an important role in achieving the superconducting state. Further, Tc attains a
value zero when atomic mass M approaches infinity. This implies that non-zero transition temperature Tc is
a consequence of finite mass of ions. According to second observation, the jump in the value of specific heat
at transition temperature indicates the presence of an energy gap in the energy specturm of the electron in the
superconducting state. We discuss below the quantum theory of superconductors due to BCS.
When an electron approaches a positive ion core then it suffers attractive Coulomb interaction. So, the ion
core sets in motion due to this attraction. Consequently the lattice gets distorted. It is obvious that this
distortion will be greater if the mass of the positive ion core is small. Now suppose that another electron
interacts with the distorted lattice. Due to this interaction, the energy of the second electron gets lowered.
Thus we can say that the two electrons interact via the lattice distortion or the phonon field. So, lowering of
energy for the electrons takes place which implies that the force between the two electrons is attractive. This
type of interaction is called as electron-lattice-electron interaction. This interaction would be strongest when
the two electrons have equal and opposite momenta and spins. The above interaction can be interpreted as the
electron electron interaction through phonons as the mediator, because the oscillatory distortion of lattice is
quantized in terms of phonons.
The superconductivity occurs when an attractive interaction (as mentioned above) between two electrons
due to phonon exchange dominates the usual repulsive interaction. This is the fundamental postulate of the
BCS theory. As discussed earlier also, the two electrons, which interact attractively, are called a Cooper
pair. The energy of the pair of electrons in bound state is less than the energy of the pair in free state. This
difference of energy is called binding of Cooper pairs. It means by applying this amount of energy we can
break this pair. The pairing is complete at T = 0 K and completely broken at T = Tc. It was observed that the
binding energy of the Cooper pair is maximum when electrons forming the pair have opposite momenta and
spins.
In addition, there is a BCS wavefunction composed of particle pairs. When treated by BCS theory, it gives the
familiar electronic superconductivity observed in metals and exhibits the energy gap. For the accomplishment
of BCS wavefunction, we need that
(i) An interaction (attraction) between electrons can lead to a ground state separated from excited
states by an energy gap. The thermal properties and most of the electromagnetic properties are
consequences of this energy gap.
(ii) The magnetic flux through a superconducting ring is quantized and effective unit of charge is 2e
rather than e.
20.11 coHErEncE lEngTH LO7

The concept of coherence is related to the idea that superconductivity is due to the mutual interaction and
correlation of the behavior of electrons that extend over a considerable distance. The coherence length is the
maximum distance up to which the states of pairs of electrons are correlated to produce superconductivity.
It is represented by Œ0. In view of this, Œ03 represents a volume known as coherence volume. The properties
of a superconductor depend on the correlation of electrons within this coherence volume. A large number of
electrons in such a volume forms Cooper pairs and acts together in the superconductivity so that a transition
is very sharp. The ratio of London penetration depth to the coherence length is given by
lL
kk =
e0
where kk is a number known as Ginzburg-Landau parameter. This number demarcates the two types of
superconductors. For type-I superconductors, 0 < kk < 1/ 2 and for type-II superconductors, kk > 1/ 2 . BCS
theory says that the coherence length Œ0 is related to the energy gap according to
hn F
e0 =
2p ◊ 2 D
Here 2D is the energy gap and vF is the Fermi velocity, given by vF = 2 EF / m .

20.12 HIgH TEMPEraTurE (Hi-Tc) SuPErconducTIvITy LO8

The high temperature superconductors represent a new class of materials, those bear extraordinary super
conducting and magnetic properties. These types of superconductors have great potential for wide ranging
technological applications. The need to keep the superconducting circuits at Liquid helium temperatures has
revealed that all such applications are expensive and specialized. Bendroz and Muller discovered a class of
superconductors with higher critical temperatures. The high temperature superconductors have transition
temperature above 40 K or may be even 90 K. Liquid nitrogen is better coolant with temperature (Tc > 77K)
than helium because of its larger heat capacity. It is inexpensive also. Since it has a temperature of 77 K, the
new materials can be maintained in their superconducting state relatively easily and cheaply. There are at
least four different crystal structures that result in high temperature superconductivity, all composed of layers
of copper-oxygen structures sandwiched around layers of other elements. These are as follows
(i) Yttrium barium copper oxide (YBa2Cu3O7) with 92 K as Tc.
(ii) Bismuth strontium calcium copper oxide {(BiPb)2Sr2Ca2Cu3Ox} with 105 K as Tc.
(iii) Thallium barium calcium copper oxide (TlBa2Ca2Cu3Oy) with 115 K as Tc.
(iv) Mercury barium calcium copper oxide (HgBa2Ca2Cu3Oy) with 135 K as Tc.
In view of their diverse applications, there has been a tremendous effort to find still more of them with even
high Tc. The hi-Tc materials are all type-II superconductors and have extremely high upper critical fields.
However, new high temperature superconductors have a drawback as they are very brittle and do not carry
enough current. Moreover, in the new hi-Tc superconductors, the formation of electrons pairs is not caused by
the same electron lattice interaction as in the BCS theory. So far it not well understood that what does cause
the formation of pairs. Intense research is being conducted in this direction.
20.13 aPPLicatiOn Of suPErcOnductivity LO8

The attractive property of zero resistance (absence of Joule heating) and Meissner effect make superconductors
very useful for many applications. Some of the applications are listed below.
(i) Superconducting cables can be used to transmit electric power over long distances without power
losses.
(ii) Superconductors can be used to check unwanted magnetic flux by using their diamagnetic property.
(iii) By using superconductor, very fast and accurate computers can be developed.
(iv) Superconductors are used to detect brain wave activities.
(v) Superconductors are also used in harnessing the various forms of nuclear energies.
(vi) Most successful application is to use superconducting wire to carry the large currents in high field
electromagnets. Such superconducting magnets are widely used in modern particle accelerators.
Superconducting magnets are also used in magnetic resonance imaging (MRI) in medicine.
(vii) The superconducting quantum interference device (SQUID) may be configured as a magnetometer
to detect incredibly small magnetic fields, i.e., small enough to measure the magnetic fields in living
organisms.
(viii) Highly powerful strong field superconductor electromagnets are fabricated using liquid helium
superconductors. These electromagnets are used in NMR spectrometers and NMR imaging that
are employed in medical diagnosis. Electromagnets are also used to produce Josephson’s devices,
electromagnetic shields and magnetically levitating world’s fastest trains.
(ix) Low temperature superconductors have been used to construct fractional wavelength antennas,
leading to a significant improvement in radiation efficiency. However, the use of liquid helium as a
cryogen limits the application of such antennas.
(x) It is seen that conventional metal guides at mm wavelength have attenuations of the order of 10
dB/m due to the high value of surface resistance of the metal walls at ~ 200 GHz. Therefore,
potential application for superconductors is also in the construction of electromagnetic waveguides.
The advantage over conventional metal waveguides would be at the higher frequencies.
(xi) Now it has been possible to design ceramic superconductors which can act at temperature > 77 K,
i.e., it can act as hi-Tc superconductors. These superconductors have advantage over low Tc
superconductors because liquid nitrogen can be used as coolant, which is cheaper and has better
cooling due to its high thermal capacity.
(xii) Other industrial applications of superconductors are through magnets, sensors, transducers and
magnetic shielding.
(xiii) Superconductors also have applications in power generation, energy storage, fusion, transformers
and transducers.
S UMMARY

✦ There are several factors that contribute to the electrical resistivity of a solid. For example, the
deviations from a perfect lattice can scatter the electrons. Also, the vibrations of lattice ions (phonons)
take place in normal modes. So, the electrons and phonon interaction scatters conduction electrons and
hence causes more resistance. Therefore, it is clear that the electrical resistance of a solid will decrease
if we cool the solid. It is amazing that electrical resistivity of some metals, alloys and compounds drops
suddenly to zero when they are cooled below a certain temperature (called critical temperature). This
phenomenon is called superconductivity.
✦ Properties of superconductors viz. electrical property, magnetic property and thermal property, were
talked about. An important property of the superconducting phase is the repulsion of magnetic flux lines
from the bulk of superconductor. It is called Meissner effect. When a specimen is placed in a magnetic
field, the magnetic flux lines pass through it. Now, if the temperature is decreased below the transition
temperature (Tc), it expels all the magnetic flux lines from inside of the specimen.
✦ Depending upon the magnetisation behaviour, the superconductors can be classified into two groups,
namely type-I or soft superconductor and type-II or hard superconductor. Type-I superconductor expels
all the magnetic field abruptly from the interior and becomes an ideal diamagnetic material. However,
type-II superconductor loses magnetisation gradually rather than abruptly.
✦ The Maxwell’s equations of electromagnetic waves were not able to explain the “zero resistance” and
“ideal diamagnetism” of superconductors. However, London and London in the year 1935 derived
two new equations to explain the superconducting state of matter. These equations are known as first
London equation and second London equation, given below respectively

dJ ne 2    ne 2 
= E and — ¥ J = - B
dt m m
✦ Effects of isotopes and magnetic field on the superconductors were discussed.
  ne 2 
✦ The second London equation — ¥ J = - B explains well the Meissner effect. Then it was proved
m
that this equation also predicts the penetration of supercurrent and magnetic flux in a superconductor.
Calculations of this penetration depth were done and it was obtained that the penetration depth is given
m
by l L = .
m0 e 2 n
✦ The generation of Cooper pair was made clear based on the interaction of electrons with the phonons.
Cooper pairs are formed due to the electron lattice interactions. The two electrons of a cooper pair have
equal and opposite momenta. They also have opposite spins, i.e., one electron is spin up and another is
spin down. Thus, the bound Cooper pair is a spin zero object and is a boson.
✦ Very late after the prediction of Einstein in 1924 that bosons could condense in unlimited numbers
into a single ground state, anomalous behaviour of liquid helium was noticed at low temperatures.
A remarkable discontinuity in heat capacity of helium was observed when it was cooled to critical
temperature of 2.17 K. The liquid density dropped and a fraction of the liquid became a superfluid with
zero viscosity. This superfluidity took place due to the fraction of helium atoms which condensed to the
lowest possible energy. This was referred to as Bose Einstein condensation.
✦ The basis of a quantum theory of superconductivity was led by the classic 1957 paper of Bardeen, Cooper
and Schrieffer, which is now called the BCS theory. The formulation of BCS theory is based on two
experimental facts viz the isotope effect and the variation of specific heat of superconductors.
✦ The superconductivity occurs when an attractive interaction between two electrons due to phonon
exchange dominates the usual repulsive interaction. This is the fundamental postulate of the BCS theory.
The energy of the pair of electrons in bound state is less than the energy of the pair in free state. This
difference of energy is called binding energy of the Cooper pairs. It means by applying this amount of
energy we can break this pair. The pairing is complete at T = 0 K and completely broken at T = Tc. It
was observed that the binding energy of the Cooper pair is maximum when electrons forming the pair
have opposite momenta and spins. In addition, there is a BCS wavefunction composed of particle pairs.
✦ The concept of coherence was introduced. The coherence length is the maximum distance up to which
the states of pairs of electrons are correlated to produce superconductivity. It is represented by e0. The
lL
ratio of London penetration depth to the coherence length, given by kk = , demarcates the two types
e0
of the superconductors. For type-I superconductors, 0 < kk < 1/ 2 and for type-II superconductors,
kk > 1/ 2 .
✦ The high temperature (hi-Tc) superconductors represent a new class of materials, those bear extraordinary
superconducting and magnetic properties. Bendroz and Muller discovered a class of superconductors
with higher critical temperatures. The hi-Tc superconductors have transition temperature above 40 K or
may be even 90 K. Liquid nitrogen is better coolant with temperature (Tc > 77K) than helium because
of its larger heat capacity. It is inexpensive also. Since it has a temperature of 77 K, the new materials
can be maintained in their superconducting state relatively easily and cheaply.
✦ Finally the applications of superconductors in medical, electronics, industry, power generation,
transportation etc. were talked about. It was felt that the superconductivity has diverse applications in
different areas of science and engineering.
S OLVED E XAMPLES
E xamplE 1 The critical temperature of lead is 7.2 K. Determine the penetration depth in lead at 5.1 K if the
penetration depth at 0 K is 380 Å.
Solution Given Tc= 7.2 K and l0 = 380 Å, l(5.1 K) = ?
Formula used is
-1/ 2
È Ê T ˆ4˘
l (T ) = l0 Í1 - Á ˜ ˙
ÍÎ Ë Tc ¯ ˙˚
-1/ 2
È Ê 5.1 ˆ 4 ˘
l (5.1 K ) = 380 Í1 - Á ˜ ˙
Î Ë 7.2 ¯ ˚
= 439.29 Å
E xamplE 2 Determine the transition temperature and critical field at 4.2 K for a given specimen of a
superconductor if the critical fields are 1.41 ¥ 105 and 4.205 ¥ 105 amp/m at 1.41 K and 12.9 K, respectively.
Solution Given HC1= 1.41 ¥ 105A/m at T1 =14.1 K and HC2 = 4.205 ¥ 105A/m at T2 = 12.9 K
Formula used is
È Ê T ˆ2˘
H c = H 0 Í1 - Á ˜ ˙
ÍÎ Ë Tc ¯ ˙˚
Thus, the critical fields at temperatures T1 and T2 can be written
È Ê T ˆ2˘ ÈT 2 - T 2 ˘
H C1 = H 0 Í1 - Á 1 ˜ ˙ = H 0 Í c 2 1 ˙
ÍÎ Ë Tc ¯ ˙˚ ÍÎ Tc ˙˚
ÈT 2 - T 2 ˘
HC 2 = H0 Í c 2 2 ˙
ÎÍ Tc ˚˙
From the above equations,we get
H C1 Tc2 - T12 1.41 ¥ 105 Tc2 - (14.1) 2

= or =
H C 2 Tc2 - T22 4.205 ¥ 105 Tc2 - (12.9) 2
Tc = 14.67 K
Using the values,

Hc = 1.41 ¥ 105 A/m, T1 = 14.1 K and Tc = 14.67 K
È Ê T ˆ2˘
H C = H 0 Í1 - Á 1 ˜ ˙
ÎÍ Ë Tc ¯ ˚˙
È Ê 14.1 ˆ 2 ˘
1.41 ¥ 105 = H 0 Í1 - Á ˜ ˙
Î Ë 14.67 ¯ ˚
or H0 = 18.504 ¥ 105 A/m
The critical field at T = 4.2 K and Tc = 14.67 K
È Ê T ˆ2˘ È Ê 4.2 ˆ 2 ˘
H c = H 0 Í1 - Á ˜ ˙ = 18.504 ¥ 105 ¥ Í1 - Á ˜ ˙
ÎÍ Ë Tc ¯ ˚˙ Î Ë 14.67 ¯ ˚
= 16.99 105 A/m
ExamplE 3 Assuming that the critical magnetic field depends upon T, find the critical current density for 1.0 mm
diameter wire of lead at 4.2 K. Take critical temperature for lead as 7.18 K and H0 for lead as 6.51 ¥ 104 A/m.
Solution Given T = 4.2 K, H0 = 6.51 ¥ 104 A/m and Tc = 7.18 K
È Ê T ˆ2˘ È Ê 4.2 ˆ 2 ˘
H c = H 0 Í1 - Á ˜ ˙ = 6.51 ¥ 104 ¥ Í1 - Á ˜ ˙
ÍÎ Ë Tc ¯ ˙˚ Î Ë 7.18 ¯ ˚
= 4.29 ¥ 104 A/ m
d 1.0
Critical current Ic = 2prHc, r = = mm = 0.5 mm
2 2
Ic 2p rH c 2 H c
Current density J c = = =
pr2 pr2 r
2 ¥ 4.29 ¥ 104
= = 1.716 108 A/m 2
0.5 ¥ 10-3
E xamplE 4 The critical temperature Tc for Hg with isotopic mass 199.5 is 4.185 K. What will be its critical
temperature when its isotopic mass is increased to 203.4.
Solution Given Tc1 = 4.185 K, M1 = 199.5 and M 2 = 203.4
Tc2 = ?
Formula used is
TcM1/2 = constant
Therefore, Tc M11/ 2 = Tc M 21/ 2

1 2
4.185 ¥ (199.5)1/ 2 = Tc2 ¥ (203.4)1/ 2
\ Tc2 = 4.1446 K
E xamplE 5 Determine the penetration depth in mercury at 0 K, if the critical temperature of mercury is 4.2 K
and the penetration depth is 57 nm at 2.9 K.
Solution Given Tc = 4.2 K
l(2.9) = 57 nm
l(0) = ?
Formula used is
1/ 2
È Ê T ˆ4˘
l (0) = l (T ) Í1 - Á ˜ ˙
ÍÎ Ë Tc ¯ ˙˚
1/ 2
È Ê T ˆ4˘
l (0) = l (2.9) Í1 - Á ˜ ˙
ÎÍ Ë Tc ¯ ˚˙
1/ 2
È Ê 2.9 ˆ 4 ˘
l (0) = 57 ¥ 10-9 ¥ Í1 - Á ˜ ˙
Î Ë 4.2 ¯ ˚
= 57 ¥ 10-9 ¥ [1 - (0.6905) 4 ]1/ 2
= 50.10 nm
E xamplE 6 Determine the critical temperature of aluminium if the penetration depth for aluminium is 16 nm
and 96 nm at 2.18 K and 8.1 K, respectively.
Solution Given l(2.18) = 16 nm and l(8.1) = 96 nm.
Formula used is
-1/ 2
È Ê T ˆ 4˘
l (T ) = l0 Í1 - Á ˜ ˙
ÎÍ Ë Tc ¯ ˚˙
96 l0 [1 - (8.1/ Tc ) 4 ]-1/ 2
=
16 l0 [1 - (2.18/ Tc ) 4 ]-1.2
4
Ê 2.18 ˆ
1- Á
Ë Tc ˜¯
6= 4
Ê 8.10 ˆ
1- Á
Ë Tc ˜¯
1
or 35 = {154968.19 - 22.58}
Tc4
Tc = 8.16 K
E xamplE 7 The critical temperature of a given superconducting sample is 1.19 K with mass 26.91. Determine
the critical temperature when the isotope mass changes to 32.13.
Solution Given Tc1 = 1.19 K, M1 = 26.91 and M 2 = 32.13, Tc2 = ?
Formula used is
Tc1 M11/ 2 = Tc2 M 21/ 2
Tc1 M11/ 2 1.19 ¥ (26.91)1/ 2
Tc2 = =
M 21/ 2 (32.13)1/ 2
1.19 ¥ 5.127 6.173
= =
5.668 5.668
Tc2 = 1.089 K
E xamplE 8 Considering the critical temperature of mercury as 4.2 K, calculate the energy gap in eV at
T = 0. Also find the wavelength of a photon whose energy is just sufficient to break up Cooper pairs in
mercury at T = 0. Find the region of the electromagnetic spectrum where such photons may be observed.
Solution Cooper pair binding energy or the energy gap is Eg = 3kTc
fi Eg = 3 ¥ 1.38 ¥ 10–23 ¥ 4.2 = 1.74 ¥ 10–22 J
1.74 ¥ 10-22
or Eg in eV = = 1.08 ¥ 10-3 eV
1.6 ¥ 10-19
The wavelength l of a photon of energy Eg
hc Ê hc ˆ
l= ÁË E g = ˜¯
Eg l
6.6 ¥ 10 -34 ¥ 3 ¥ 108
= = 1.14 ¥ 10 -3 m
1.74 ¥ 10 -22
From the value of l, it is clear that these photons are in the very short wavelength part of the microwave region.
E xamplE 9 In Example 8, does the metal look like a superconductor to electromagnetic waves having
wavelengths shorter than 1.14 ¥ 10–3 m?
Solution Since the energy content of shorter wavelength photons is more than sufficient to break up the Cooper pairs or
to excite the conduction electrons through the energy gap into the non-superconducting states above the gap, the metal
would not work as superconductor to the said electromagnetic waves.
Q.1 The phenomenon of superconductivity was discovered by Kamerlingh Onnes in

(a) 1931 (b) 1911 (c) 1921 (d) 1811
Q.2 The temperature at which a normal material turns into a superconducting state is called the
(a) critical temperature (b) superconducting temperature
(c) high temperature (d) low temperature
Q.3 The transition temperature Tc of all superconducting materials satisfies the relation
(a) T < Tc (b) T > Tc (c) T = Tc (d) none of these
Q.4 Resistance of a substance in superconducting state becomes
(a) infinite (b) finite (c) zero (d) arbitrary
Q.5 When a superconducting material is placed in an external magnetic field, it

(a) repels the magnetic field lines (b) does not influence magnetic field lines
(c) attracts the magnetic field lines (d) enhances the magnetic field lines
Q.6 In superconducting state, the material becomes
(a) ferromagnetic (b) perfect diamagnetic
(c) strong paramagnetic (d) ferrimagnetic
Q.7 If the bulk specimen passes through the normal state to the superconducting state, its entropy
(a) increases (b) decreases
(c) remains unaffected (d) becomes half of it
Q.8 The variation of critical temperature with increasing atomic mass is called
(a) isotope effect (b) isobar effect (c) isomeric effect (d) none of these
Q.9 State of superconducting is generally achieved at
(a) very low temperatures (b) low temperatures
(c) high temperature (d) superhigh temperatures
Q.10 The London penetration depth is given by
1/ 2 1/ 2
Ê m ˆ Ê m ˆ
(a) l = (b) l = (nmm0e2)1/2 (c) l = (d) none of these
ÁË nm e 2 ˜¯ ÁË mm e 2 ˜¯
0 0
Q.11 Superconductor is
(a) an ideal conductor (b) an ideal diamagnetic
(c) not an ideal diamagnetic (d) both (a) and (b)
Q.12 The magnetic susceptibility of a superconductor becomes
(a) positive (b) zero (c) negative (d) none of these
Q.13 Which of the following statement is correct interms of critical field strength of superconductor?
(a) It depends on mechanical stress (b) It varies with temperature
(c) It depends on the purity of metal (d) All of these
Q.14 Which of the following relation represents London’s equation?
 
 ∂H  H
(a) curl E = - m0 (b) curl J s = - 2
∂t lL
2
 l
(c) curl J s = - L (d) both (a) and (b)
H
Q.15 Which of the following relation is correct?
È Ê T ˆ2˘ È Ê T ˆ2˘
(a) H c = H 0 Í1 + Á ˜ ˙ (b) H c = H 0 Í1 - Á ˜ ˙
ÍÎ Ë Tc ¯ ˙˚ ÎÍ Ë Tc ¯ ˚˙
È Ê T ˆ˘
(c) H c = H 0 Í1 - Á ˜ ˙ (d) None of these
Î Ë Tc ¯ ˚
Q.16 Which of the following relation is correct?
 ne 2 
(a) Cn(T) = gT + bT3 (b) curl J = - B
-1/2
m
l (T ) È Ê T ˆ ˘
4
(c) = Í1 - Á ˜ ˙ (d) all of these
l (0) ÍÎ Ë Tc ¯ ˚˙
Q.1 What do you understand by superconductivity?

Q.2 Explain whether a superconductor is a perfect diamagnet.
Q.3 What is Meissner effect?
Q.4 Mention some important changes that occur in substances when they change from normal state to
superconducting state.
Q.5 Name the types of superconductors?
Q.6 What is isotopic effect in superconductors?
Q.7 What are high Tc superconductors? Give some examples.
Q.8 Discuss how Cooper pairs are formed?
Q.9 What do you understand by phonons?
Q.10 What is the binding energy of a Cooper pair?
P RACTICE P ROBLEMS
general Questions
Q.1 What is superconductivity? Give the main properties of a superconductor?
Q.2 What do you mean by critical field in superconductivity?
Q.3 Explain the distinction between type-I (soft) and type-II (hard) superconductors.
Q.4 Explain the difference between the type-I and type-II superconductors using Meissner effect.
Q.5 What do you mean by superconductivity? Describe the effect of (a) magnetic field (b) frequency
(c) isotopes on superconductors.
Q.6 Define and explain the Meissner effect in superconductors.
Q.7 Describe the effect of an external magnetic field on the superconducting state of a material. What do
you mean by flux exclusion and what is Meissner effect?
Q.8 What is the significance of critical temperature, critical magnetic field and critical current density for
superconductors?
Q.9 Derive the London equations and discuss how they explain Meissner effect and flux penetration?
Q.10 Give brief outline of Bardeen, Cooper and Schrieffer (BCS) theory of superconductivity. Show that
this theory provides adequate explanation of superconducting state.
Q.11 Describe major uses and potentialities of superconductors.
Q.12 Why superconductivity is a low temperature phenomenon?
Q.13 What do you mean by high temperature superconductivity?
(i) Penetration of magnetic field in a superconductor and penetration depth.
(ii) Flux quantization.
Q.15 Discuss Bose Einstein condensation.
Q.16 Explain why the electrical resistivity of a solid goes down with the increase in temperature?
X-Rays 21
Learning Objectives
LO 1 Know about origin of X-rays L0 3 Explain Moseley’s Law on the basis of
L0 2 Learn about properties of X-rays, Bohr's theory and its importance
continuous X-ray spectrum, line L0 4 Discuss practical applications of X-rays
spectrum
Introduction
In the phenomenon of photoelectric effect, a photon can eject an electron from a metallic surface if
its energy is greater than the threshold value. As is clear from this process, the photons of light can
transfer energy to the electrons. However, in 1895 Roentgen observed the reverse process in which the
kinetic energy of moving electrons was converted into photons under suitable conditions. He observed
that when fast electrons impinge on the anode material in the Crooks discharge tube, some rays are
produced that have highly penetrating power. These rays were named as X-rays. Actually in this process,
the electrons passing near a nucleus in the target are decelerated and hence emit a continuous spectrum
of radiation (Bremsstarhlung) ranging up from a minimum wavelength. In addition to this, the electrons
may eject an electron from an inner shell of a target atom. Then the resulting transition of an electron of a
higher energy level to this level produces radiation of specific wavelength, which is the characteristic X-ray
spectrum of the target and is specific to the target element.
These rays are found to penetrate paper, thick wooden blocks, glass and thin metal sheet, etc. The
penetrating power of these rays depends on the speed of the moving electrons, i.e., faster the moving
electrons the greater the penetration power of the X-rays. The intensity of X-ray beam is found to be in
direct proportion to the number of electrons, i.e., the intensity of X-ray beam is greater if the number of
electrons is larger.
21.1 Origin Of X-rays LO1

X-rays are produced when fast moving electrons strike with H.T.
a target of high melting point and high atomic number, for
example, tungsten, platinum, molybdenum, copper and
chromium. A Coolidge tube is generally used to produce C T Inlet
X-rays, as shown in Fig. 21.1. It is evacuated at low pressure
of 10–5 mm of Hg. This tube contains a cathode C and target Outlet
T. The filament is heated with a low voltage source to emit
L.T.
electrons. A high voltage (ª 10,000 – 30,000 volts) is applied X-rays
between the filament and target to accelerate the electrons.
When these high velocity electrons strike the target, X-rays
are produced. The wavelength of X-rays produced depends Figure 21.1
on various factors like nature of target material, and initial and final energies of incident electrons.
Only about 0.2% of the electron beam energy is used to produce X-rays, the rest of the energy transforms into
heat. In order to remove the heat so produced by the beam, the target material is mounted on a hollow copper
tube through which the cold water is flowing continuously.
21.1.1 Control of intensity
The intensity of X-rays is controlled by controlling the intensity of incident electrons striking the target. As
mentioned earlier, actually greater number of the incident electrons produce more intense X-rays. In other
words, the intensity of X-rays can be controlled by controlling the current in the filament.
21.1.2 Control of Penetrating Power
The penetrating power of X-rays is controlled by controlling the potential difference between the filament and
target. If the potential difference (V) is increased, the energy of incident electrons increases. This results in the
X-rays of higher energy and hence their higher penetration power.
21.2 PrOPerties Of X-rays LO2

Properties of X-rays are listed below.
(i) X-rays are electromagnetic radiations of very short wavelength of the order of 10–10 metre and they
travel with the velocity of light.
(ii) They are not deflected in the presence of electric and magnetic fields. It means they are not charged
particles like electrons.
(iii) They produce fluorescence in some substances like barium platino-cyanide, zinc sulphide, etc.
(iv) They affect the photographic plate and the effect is more intense than light.
(v) X-rays ionize the gas when they are passing through it.
(vi) They produce photoelectric effect.
(vii) They can pass through many solids due to high penetrating power.
(viii) They show interference, diffraction and polarisation like light.
X-Rays 737
21.3 X-ray sPeCtra LO2

X-ray spectra consist of two types of spectra, namely continuous spectrum and characteristic or line spectrum.
21.3.1 Continuous X-ray spectrum
The continuous X-ray spectrum can be produced at sufficiently low potential Electron
applied between the target and filament. Some of the fast moving electrons mn2/2
penetrate the core of the atoms of the target material, as shown in Fig. 21.2.
During this process, the electrons get decelerated and their velocities are
reduced. According to the classical theory, the decelerated charge particle
radiates energy in the form of X-rays. Due to different losses in the energy of mn¢2/2
incident electrons, the radiations of all possible wavelengths within a certain
range are emitted which constitute a continuous spectrum. Figure 21.2
If the velocity of the incident electron reduces from v to v¢, then the frequency (n) of the emitted X-ray photon
is given by
1 2 1
mv - mv ¢ 2 = hn (i)
2 2
If the incident electron is completely stopped, then the frequency of the emitted X-ray photon will be
maximum, i.e.,
1 2
mv = hn max = eV (ii)
2
where, V is accelerating potential of electron,
hc Ê c ˆ
or = eV ÁË∵ n max = l ˜¯
lmin min
hc
or lmin =
eV (iii)
Eq. (iii) gives the minimum wavelength limit of the continuous X-ray spectrum. This is also called quantum
limit. Putting the values of h, c and e in the equation, we get
6.62 ¥ 10 -34 ¥ 3 ¥ 108
lmin =
1.6 ¥ 10 -19 ¥ V
12400
lmin = (iv)
V
Thus, Eq. (iv) shows that lmin is inversely proportional to the voltage (V) applied between the cathode and
target. If V is in volts, lmin is obtained in Å.
30 kV
21.3.1.1 Features of Continuous X-rays Spectrum
20 kV
Some features of continuous X-ray spectrum are given below.
(i) X-rays are produced due to the deceleration of fast moving
Intensity(I)
10 kV
electrons.
(ii) The intensity of continuous spectrum increases as the
potential is increased (Fig. 21.3). lmin lmin lmin
Wavelength l
(iii) The minimum wavelength limit shifts towards lower
wavelength as the potential is increased. Figure 21.3
21.3.2 Characteristic X-ray spectrum: Line spectrum

It consists of some well-defined wavelengths superimposed on the continuous spectrum. The spectrum lines
generally occur in the form of small groups. The characteristic X-ray spectrum is produced when extremely high
energetic electrons penetrate well inside the atoms of the target and collide with tightly bound electrons of innermost
K orbit of the atom, as shown in Fig. 21.4. If this vacancy is filled up by the electron of the second orbit L, Ka line
is produced and if it is filled up by the electron of the third orbit M, Kb line is produced and so on. In this manner
we get different series like K, L, M etc. in the emission of X-ray spectra, as shown in Fig. 21.5. Therefore, the
energy of the emitted X-ray photon corresponding to Ka line is given by
hn Ka = EK - EL (v)
electron
Hole X-ray
Hole
Photon
hn
+ +
K electron
L
M
(a) (b)
where, EK and EL are the energies of electrons required to remove Characteristic

K
electrons from K and L orbits, respectively. Fig. 21.6 depicts the X-ray
plot of intensity I versus wavelength l for characteristic X-ray L
spectrum superimposed on the continuous spectra.
M
Intensity (I)
21.3.2.1 Features of Characteristic X-ray Spectrum

Some features of characteristic X-ray spectrum are given below
(1) Line spectrum is produced when electrons of innermost
orbit are removed and this is followed by jumping of
electrons from outer orbits. Wavelength l
(2) The wavelength of the lines of K-series shifts towards Figure 21.6
lower values as the atomic number of the target element increases.
(3) The intensity of definite spectral line depends on the probability of that particular transition.
21.4 MOseLey’s Law LO3

Moseley analysed characteristic X-ray spectra emitted by targets of a number of heavy elements (Z = 22 to
Z = 30) and established a simple relation between them. He observed that the characteristic X-ray spectra
of different elements consist of K, L, M series. Also, he found that on increasing the atomic number Z, the
characteristic X-ray spectral lines shift towards shorter wavelengths side or towards higher frequencies side.
X-Rays 739
Moseley plotted a graph of the K-series lines of the characteristic X-ray spectra for a number of elements
between n (n is the frequency) and Z of the target material and found almost a straight line (Fig. 21.7).
From this he concluded that
Kb
n µ Z or n µ Z 2
This conclusion is known as Moseley’s law. Mathematically √n
Ka
n = a(Z – b)2
where, a and b are the constants for the given transition of the K-series.
The constant b is known as screening constant.
21.4.1 explanation based on Bohr’s theory
According to the Bohr’s theory, b Z
Figure 21.7
c Ï1 1¸
n = = RcZ 2 Ì 2 - 2 ˝ ,
l Ó n1 n2 ˛
where, R is Rydberg constant and Z is the atomic number of the element. When an electron jumps from
L-level to K-level, i.e., n2 = 2 and n1 = 1, the frequency of Ka line is given by
Ï1 1¸ 3
n = RcZ 2 Ì 2 - 2 ˝ = RcZ 2
Ó1 2 ˛ 4
Now, if one of the two electrons of the K-shell of an atom is knocked off, an electron from L-shell would make
a transition to the K-shell, thereby emitting Ka line. However, the remaining K-electrons produce a screening
effect which reduces the force of attraction of the nucleus for the L-electron. Therefore, for L-electron the
effective charge of the nucleus reduces to (Z – l)e. Hence, replacing Z by Z – 1 in Bohr’s formula, we get
3
n= Rc( Z - 1) 2
4
For spectral line La, n2 = 3, n1 = 2
Ê 1 1ˆ
n = Rc( Z - 7.4) 2 Á 2 - 2 ˜
Ë2 3 ¯
5
= Rc( Z - 7.4) 2
36
In general, we can write
È1 1˘
n = Rc( Z - b) 2 Í 2 - 2 ˙
Î n1 n2 ˚
or n = a ( Z - b) 2
or n µ ( Z - b)
21.4.2 importance of Moseley’s Law

The importance of Moseley’s law is listed below.
(a) It yields a unique method for identifying as well as for classifying different elements in the periodic
table. According to Moseley’s law, the atomic number Z of the element determines the physical and
chemical properties of the element. Therefore, the elements in the periodic table must be arranged
according to Z instead of atomic weight. Infact before Moseley, Mendeleef constructed periodic
table by arranging elements in the order of increasing atomic weight and he put
58.69
28Ni before 27Co58.94, 19K39 before 18A40
and this anomaly was removed by Moseley by putting them according to their atomic number.
(b) It is helpful in determining Z of rare earth elements and position in the periodic table.
(c) Based on this new elements like rhenium (75), hafnium (72), promethium (61), technetium (43) etc.
were discovered.
21.5 PraCtiCaL aPPLiCatiOns Of X-rays LO4

We mention below some of the important applications of X-rays.
(i) The X-rays are extensively used in surgery to detect fractures, diseased organ, etc. as these rays can
pass through flesh and not through bones.
(ii) Since the X-rays kill the diseased tissues of the body, they are used to cure some skin diseases,
malignant tumors, etc.
(iii) The X-rays are used to detect defects in motor tyres, golf and tennis balls, etc. and also for testing
the uniformity of insulating materials.
(iv) People have explored the high penetrating power of the X-rays to investigate the structure of metals,
crystalline solids, alloys, complex organic molecules, etc.
(v) The X-rays are used for the detection of explosive arms at customs posts.
(vi) X-ray spectrography and X-ray crystallography are the most recent applications of the X-rays.
s UmmarY

✦ For X-ray generation, a process is used in which the kinetic energy of moving electrons is converted into
photons under suitable conditions.
✦ Bremsstarhlung arises from the retardation of charged particles, commonly electrons, stopped by the
target of high atomic number.
✦ The X-rays can pass through the opaque materials and cause phosphorescent substances to glow.
✦ The penetrating power of the X-rays depends on the speed of the moving electrons, i.e., faster the
moving electrons greater the penetration power of the X-rays.
✦ The intensity of X-ray beam is found to be in direct proportion to the number of electrons, i.e., the
intensity of X-ray beam is greater if the number of electrons is larger.
✦ The continuous X-ray spectrum is produced at sufficiently low potential applied between the target and
filament of the Coolidge tube. Here, some of the fast moving electrons penetrate the core of the atoms
of the target material, get decelerated and radiate energy in the form of X-rays.
X-Rays 741
✦ The minimum wavelength limit of the continuous spectrum, which is also called quantum limit, is
hc
given by lmin = , where V is the accelerating potential of the electron.
eV
✦ The characteristic X-ray spectrum is produced when extremely high energetic electrons penetrate well
inside the atoms of the target and collide with tightly bound electrons of innermost K orbit of the atom.
If this vacancy is filled up by the electron of the second orbit L, Ka line is produced. If this vacancy is
filled up electron of the third orbit M, Kb line is produced. Similarly La line is produced if the vacancy
in L orbit is filled up by the electron of the third orbit M.
✦ Moseley analysed characteristic X-ray spectra emitted by targets of heavy elements whose atomic
numbers ranged from 22 to 30. He achieved a conclusion that satisfies the expression n = a ( Z - b).
Here a and b are the constants for the given transition of the K series.
✦ Moseley’s law was explained based on Bohr’s theory that related the frequency of radiation with the
atomic number of the element and the transitions between two energy levels.
✦ Moseley’s law had removed the anomaly took place due to the arrangement of elements in periodic table
as per the order of increasing atomic weight done by Mendeleef. Moseley suggested this arrangement
based on the atomic number rather than the atomic weight.
s olved e XamPles
E xamplE 1 An X-ray tube operates at the voltage (i) 40 kV, (ii) 20 kV, and (iii) 100 kV. Find the maximum
speed of electrons striking the anti-cathode and shortest wavelength of X-rays produced.
Solution Given (i) V = 40 kV (ii) V = 20 kV and (iii) 100 kV. Formula used is
1 2 12400
mv = eV and lmin =
2 V
2eV 2 ¥ 1.6 ¥ 10 -19
or v= = V
m 9.1 ¥ 10 -31
v = 0.593 ¥ 106 ( V ) m/sec
(i) V = 40 ¥ 103 V
12400
v = 0.593 ¥ 106 4 ¥ 104 , lmin =
4 ¥ 104
= 1.186 ¥ 108 m / sec = 0.31Å
(ii) V = 20 ¥ 103 V
12400
v = 0.593 ¥ 106 2 ¥ 104 , lmin =
2.0 ¥ 104
= 8.39 ¥ 107 m / sec = 0.62Å
(iii) V = 100 ¥ 103 V = 10 ¥ 104 V
12400
v = 0.593 ¥ 106 10 ¥ 104 , lmin =
105
= 1.875 ¥ 108 m / sec = 0.124Å
E xamplE 2 The short wavelength limit of the continuous X-ray spectrum emitted by an X-ray tube operated
at 30 kV is 0.414 Å. Calculate the Planck’s constant.
Solution Given V = 3.0 ¥ 104 V, lmin = 0.414 Å
e = 1.6 ¥ 10–19 C and c = 3 ¥ 108 m/sec
Formula used is
hc eV
lmin = or h = lmin
eV c
1.6 ¥ 10-19 ¥ 3.0 ¥ 104 ¥ 0.414 ¥ 10-10
h=
3 ¥ 108
34
6.624 10 J sec
E xamplE 3 An X-ray tube is operated at 25 kV. Calculate the minimum wavelength of X-rays emitted from
it.
Solution Given V = 25 ¥103 V
Formula used is
12400 12400
lmin = =
V 25 ¥ 103
= 0.496 Å
E xamplE 4 An X-ray tube operates at 13.6 kV. Find the maximum speed of electron striking the target.
Solution Given V = 13.6 ¥ 103 V
Formula used for maximum kinetic energy is
1 2
mv = eV
2
2eV
v= = (0.593 ¥ 106 ) V m/sec
m
= 0.593 ¥ 106 ¥ 13.6 ¥ 103
= 6.92 ¥ 107 m/sec
E xamplE 5 If the potential difference applied across an X-ray tube is 10 kV and current through it is 2.0 mA.
Calculate the velocity of electrons at which they strike the target.
Solution Given V = 10 ¥ 103 V and I = 2 ¥ 10–3 A
2eV
Formula used is v = = 0.593 ¥ 106 ¥ 10 ¥ 103
m
= 5.93 ¥ 107 m/sec
E xamplE 6 Electrons are accelerated in a television tube through a potential difference of 9.8 kV. Find the
highest frequency and minimum wavelength of the electromagnetic waves emitted, when these
strike on the screen of the tube. In which region of the spectrum will these waves lie?
X-Rays 743
Solution Given V = 9.8 ¥ 103 V

Formula used is
hc 12400
lmin = =
eV V
12400
=
9.8 ¥ 103
= 1.26 Å
c 3 ¥ 108
and n max = =
lmin 1.26 ¥ 10-10
= 2.38 ¥ 1018 Hz
E xamplE 7 If the potential difference applied across an X-ray tube is 12.4 kV and current through it is 2 mA,
calculate (i) the number of electrons striking the target per second and (ii) the speed with which they strike it.
Solution Given V = 12.4 ¥ 103 V and I = 2 ¥ 10–3 A
I 2.0 ¥ 10 -3
(i) For current I = ne or n = = or n = 1.25 ¥ 1016 electrons / sec
e 1.6 ¥ 10 -19
2 eV
(ii) n = = 0.593 ¥ 106 V m/sec
m
= 0.593 ¥ 106 ¥ 12.4 ¥ 103
= 6.6 ¥ 107 m/sec
E xamplE 8 An X-ray tube is operated at an anode potential of 10 kV and anode current of 15 mA. Calculate
(i) number of electrons hitting the anode per second and (ii) the minimum wavelength produced by the X-ray tube.
Solution Given V = 10 ¥ 103 V and I = 1.5 ¥ 10–3 A.
I 15 ¥ 10 -3
(i) For current, I = ne or n = =
e 1.6 ¥ 10 -19
= 9.38 ¥ 1016 electrons / sec
(ii) lmin = 12400 Å = 124003 Å

V 10 ¥ 10
= 1.24 Å
E xamplE 9 An X-ray tube is operated at 50 kV and current through the tube is 1.0 mA. What is the number
of electrons striking the target per second?
Solution Given I = 1.0 ¥ 10–3 A and V = 50 ¥ 103 V.
I
(i) For current I = ne or n =
e
10-3
or n=
1.6 ¥ 10-19
6.25 1015 electrons/sec
E xamplE 10 What voltage must be applied to an X-ray tube for it to emit X-rays with minimum wavelength
of (i) 40 pm and (ii) 1.0 Å
Solution Given wavelengths (i) 40 × 10–12 m (ii) 1.0 ×10–10 m
12400 12400
Formula used is lmin = Å or V =
V lmin
(i) lmin = 4.0 ¥ 10-11 m
12400 ¥ 10 -10
V= = 31 kV
4 ¥ 10 -11
(ii) lmin = 1.0 Å ¥ 10-10 m

12400 ¥ 10 -10
V= = 12.4 kV
10 -10
E xamplE 11 An X-ray tube operating at (i) 44 kV (ii) 50 kV emits a continuous spectrum with shortest
wavelength (i) 0.284 Å and (ii) 0.248 Å, respectively. Calculate the Planck’s constant.
Solution Given (i) V = 44 ¥ 103 V, l = 0.284 ¥ 10–10 m and (ii) V = 50 ¥ 103 V, l = 0.284 ¥ 10–10 m
Formula used is
hc eV
lmin = or h = ¥ lmin
eV c
1.6 ¥ 10-19 ¥ 44 ¥ 103 ¥ 0.284 ¥ 10-10

(i) h =
3 ¥ 108
= 6.66 ¥ 10-34 J sec
1.6 ¥ 10-19 ¥ 50 ¥ 103 ¥ 0.248 ¥ 10-10

(ii) h =
3 ¥ 108
= 6.61 ¥ 10-34 J sec
E xamplE 12 The K-absorption limit for Uranium is 0.1 Å. What is the excitation potential of the tube to give
this radiation?
Solution Given K-absorption limit, means lmin = 1.0 ¥ 10–11 m
12400 ¥ 10 -10
Formula used is lmin = m
V
12400 ¥ 10-10 12400 ¥ 10-10
or V= = = 124 kV
lmin 10 –11
E xamplE 13 Given that K-absorption edge for lead is 0.14Å and the minimum voltage required for producing
K-lines in lead is 88.6 kV. Determine the ratio of h/e.
Solution Given K-absorption edge for lead, lmin = 0.14 ¥ 10–10 m and V = 88.6 ¥ 103 V
Formula used is
hc h V
lmin = or = lmin
eV e c
X-Rays 745
h 88.6 ¥ 103
or = ¥ 0.14 ¥ 10-10
e 3 ¥ 108
= 4.134 ¥ 10-15 Jsec / C
E xamplE 14 Calculate the wavelength of Ka line for an atom of atomic number 92 by using Moseley’s law
and considering Rydberg constant as 1.1 ¥ 105cm–1
Solution Given Z = 92, R = 1.1 ¥ 105 cm–1
Formula used is
È1 1˘
n = cR ( Z - b) 2 Í 2 - 2 ˙
n
Î 1 n2˚
Ka line is obtained when an electron jumps from L-shell (n2=2) to K-shell (n1=1). Now b = 1 for K-series. The wavelength
l is given by
1 È1 1˘
= R ( Z - b) 2 Í 2 - 2 ˙
l n
Î 1 n2˚
È1 1˘ 3 Ê cˆ
= 1.1 ¥ 105 ¥ (92 - 1)2 ¥ Í 2 - 2 ˙ = 1.1 ¥ 105 ¥ (91)2 ¥ ÁË∵ l = ˜¯
Î1 2 ˚ 4 n
4
or l=
1.1 ¥ 3 ¥ (91)2 ¥ 105
= 1.464 ¥ 10-9 cm
= 0.15 Å
E xamplE 15 If the Ka radiation of Mo (Z = 42) has a wavelength of 0.71Å, determine the wavelength of the
corresponding radiation of Cu (Z = 29).
Solution Given lMo = 0.71 ¥ 10–10 m corresponding to Z = 42. lCu = ? corresponding to Z = 29
Formula used is
1
n= = a ( Z - b) 2 For Ka line b = 1
l
l 1
= a (42 - 1) 2 and = a (29 - 1) 2
lMo lCu
lCu a (42 - 1) 2 1681

= =
lMo a (29 - 1) 2 784
lCu = 1.52 ¥ 10-10 m

= 1.52 Å
E xamplE 16 Determine the wavelength of Ka X-rays emitted by an element having Z = 79, b = 1 and
a = 2.468 × 1015 sec–1
Solution Given Z = 79, b = 1, a = 2.468 ¥1015 sec–1
Formula used is
n = a ( Z - b) 2
n = 2.468 ¥ 1015 ¥ (79 - 1) 2
= 1.502 ¥ 1019 sec -1
c 3 ¥ 108
l= = = 1.997 ¥ 10 -11 m
n 1.502 ¥ 1019
= 0.1997 Å
E xamplE 17 Calculate the ionisation potential of K-shell electron of copper. Given that Z for copper is 29
and Rydberg’s constant for hydrogen R = 1.097 ¥ 107 m–1.
Solution Given Z = 29 and R = 1.097 ¥ 107 m–1
The frequency of Ka X-ray spectral line is given by Moseley’s law
3
v= Rc( Z - 1) 2
4
3
= ¥ 1.097 ¥ 107 ¥ 3 ¥ 108 ¥ (29 - 1) 2
4
3
= ¥ 1.097 ¥ 107 ¥ 3 ¥ 108 ¥ (28) 2
4
= 1.936 ¥ 1018 Hz
Energy of photon corresponding Ka line

6.62 ¥ 10 -34 ¥ 1.936 ¥ 1018
E = hn = eV
1.6 ¥ 10 -19
= 8.01 keV
Now
E = EK – EL, where EK is the ionisation energy, i.e., energy of electrons required to remove electrons from K-orbit
\ EK = E + EL
= 8.01 + 0.931 keV (Taking EL = 0.931 keV)
= 8.942 keV
E xamplE 18 Calculate the frequency of Ka line, when atomic number of the anti-cathode is 79. Given
R = 1.097 ¥ 107 m–1.
Solution Given Z = 79 and R = 1.097 ¥ 107m–1
Formula used is
3
n= Rc( Z - 1) 2
4
3
= ¥ 1.097 ¥ 107 ¥ 3 ¥ 108 ¥ (79 - 1) 2
4
= 1.502 ¥ 1019 s -1
X-Rays 747
E xamplE 19 Find the energy in eV and wavelength of Ka X-ray line of Co (Z = 27).

Solution Given Z = 27
Formula used for Ka line is
3
n= Rc( Z - 1) 2
4
3
= ¥ 1.097 ¥ 107 ¥ 3 ¥ 108 ¥ (27 - 1) 2
4
= 1.67 ¥ 1018 Hz
6.62 ¥ 10-34 ¥ 1.67 ¥ 1018

Energy E = hn = eV
1.6 ¥ 10-19
= 6.91 keV
Now wavelength
c 3 ¥ 108
l= =
n 1.67 ¥ 1018
= 1.79 Å
Q.1 X-rays are

(a) electromagnetic radiations (b) a stream of protons
(c) a stream of electrons (d) a stream of neutrons
Q.2 The speed of X-rays in vacuum is
(a) greater than that of the visible light (b) equal to that of the visible light
(c) smaller than that of the visible light (d) none of these
Q.3 In order to produce X-rays, an element of high atomic number is bombarded by
(a) photons (b) electrons (c) neutrons (d) protons
Q.4 Which one of the following has the lowest wavelength
(a) visible light (b) X-rays (c) infra-red rays (d) ultraviolet light
Q.5 What would be the energy of the X-ray photon when an X-ray tube is operated at 30 kV?
(a) 60 keV (b) 30 keV (c) 90 keV (d) 30 MeV
Q.6 The minimum wavelength limit of the continuous X-ray spectrum will be doubled if
(a) accelerating potential of electron is doubled.
(b) accelerating potential of electron is halved
(c) accelerating potential of electron is increased 4 times.
(d) accelerating potential of electron is reduced 1/4 times
Q.7 The minimum frequency limit of X-ray spectrum depends on
(a) nature of target material (b) degree of vacuum in Coolidge tube
(c) potential energy of incident electron (d) kinetic energy of incident electron
Q.8 In characteristic X-ray spectrum, Ka line is produced when the vacancy in K-orbit is filled by the
electron of
(a) third orbit M (b) second orbit L (c) fourth orbit N (d) none of the above
Q.9 In continuous X-ray spectrum the intensity
(a) increases as the potential is increased
(b) decreases as the potential is increased
(c) increases if the number of fast moving electrons is increased
(d) increases if the number of fast moving electrons is decreased
Q.10 The frequency of any line in characteristic X-ray spectrum is directly proportional to
(a) square root of the atomic number of target element
(b) atomic number of target element
(c) square of atomic number of target element
(d) thickness of the target
Q11. X-rays were discovered by
(a) Roentgen (b) Curie (c) Bohr (d) none of these
Q12. Penetrating power of X-rays can be increased by
(a) increasing the potential between the anode and cathode
(b) increasing the cathode filament current
(d) none of these
Q13. The wavelength of X-rays is the order of
(a) 1 Å (b) 1 m (c) 1 mm (d) none of these
Q14. Which of the following relation is correct?
hc È1 1˘
(a) lmin = (b) n = Rcz2 Í 2 - 2 ˙
eV Î n1 n2 ˚
(c) n = a(z – b)2 (d) all of these
Q15. The intensity of X-rays in Coolidge tube increases with
(a) increasing filament current
(b) increasing the potential difference between anode and cathode
(c) decreasing filament current
(d) none of these
Q16. The Bragg’s law of diffraction is
2d
(a) 2d sin q = n l/2 (b) 2d sin q = nl (c) = nl (d) none of these
sin q
Q.1 What are X-rays?

Q.2 How are the X-rays produced?
Q.3 How do you control the intensity of X-rays?
Q.4 Differentiate between characteristics X-ray spectrum and continuous X-ray spectrum.
X-Rays 749
Q.5 How will it affect the cut off wavelength of X-rays if separation between the cathode and target is
doubled?
Q.6 What is Bremsstarhlung?
Q.7 Explain the difference in origin of X-rays and visible light.
Q.8 Is it appropriate to regard X-ray production as the inverse of photoelectric effect?
Q.9 What is the difference between optical spectra and X-ray spectra?
P ractice P roblems
general Questions
Q.1 Discuss X-rays in view of their production and properties.
Q.2 Describe the construction and working of a Coolidge tube. How can you control (i) the intensity
(ii) the quality of X-rays? What are hard and soft X-rays?
Q.3 Why should anti-cathode have high atomic number and high melting point?
Q.4 What are continuous and characteristics X-rays and how are they produced? What is the minimum
wavelength limit and how it is related with the voltage applied across the X-ray tube?
Q.5 (a) Discuss the origin and mechanism of production of the continuous X-ray spectra. What is the
source of energy of photon of continuous X-rays? Show that the lowest wavelength limit of
continuos X-ray spectra is inversely proportional to accelerating potential of X-ray tube.
(b) Draw the graph of relative intensity of continuous spectra versus wavelength of X-rays and show
1
that lmin is proportional to .
V
Q.6 The potential difference between the cathode and anode in X-ray tube is doubled and the separation
between the cathode and target is also doubled. How will it affect the cut-off wavelength?
Q.7 Distinguish between continuous and characteristic X-ray spectra. Why is the characteristic spectra so
called? How is the production of characteristic X-ray spectra accounted for? Discuss the transition for
K and L series.
Q.8 What is Moseley’s law? How can it be explained on the basis of Bohr’s theory? What is its importance?
Q.9 (a) Describe Moseley’s work on X-rays. State and explain Moseley’s law. Show it graphically.
(b) Derive Moseley’s law on the basis of Bohr’s theory.
(c) Discuss the importance of Moseley’s observations of X-ray spectra of different elements. What
conclusions were drawn by him?
Q.1 An X-ray tube operates at (i) 50 kV and (ii) 18 kV. Compute the shortest wavelength of X-rays produced
and also find the maximum speed of electrons striking the target.
[Ans: (i) 0.248 Å, 1.33 ¥ 108 m/sec (ii) 0.69 Å, 7.996 ¥ 107 m/sec]
Q.2 Calculate the minimum wavelength when the potential difference applied to the X-ray tube is 98 kV.
[Ans: 0.125 Å]
Q.3 What is the shortest wavelength of X-ray produced in a tube when the applied voltage is 12.4kV?
[Ans: 1.0 Å]
Q.4 An X-ray tube operated at 40 kV emits a continuous X-ray spectrum with a short wavelength limit
lmin = 0.31Å. Calculate the Planck’s constant. [Ans: 6.61 ¥ 10–34 J sec]
Q.5 What voltage must be applied to an X-ray tube for it to emit X-rays with minimum wavelength of
(i) 0.5 Å and (ii) 0.25 Å. [Ans: (i) 24.84 kV (ii) 49.68 kV]
Q.6 An X-ray tube is operated at an anode potential 12.4 kV and current 15mA. Calculate (i) minimum
wavelength produced by X-ray tube and (ii) number of electrons hitting the anode per second.
[Ans: (i) 1.0 Å (ii) 9.4 ¥ 1016 electrons/sec]
Q.7 An X-ray tube operating at (i) 30 kV and (ii) 200 kV emits a continuous spectrum with shortest
wavelength (i) 0.414 Å and (ii) 6.2 ¥ 10–12 m. Calculate the Planck’s constant.
[Ans: (i) 6.62 ¥ 10–34 sec (ii) 6.61 ¥ 10–34 sec]
Q.8 For a platinum (Z = 78), the wavelength of La line is 1.32 Å. For an unknown element, the wavelength
of La line is 4.174 Å. Determine the atomic number of unknown element. Take b = 7.4 for La line.
[Ans: 47.1]
Q.9 The wavelength of Ka line in copper is 1.54 Å. Calculate the ionisation potential of K-shell electron in
copper. [Given the energy of L-shell as 0.923 keV.] [Ans: 8.99 keV]
Q.10 Calculate the frequency of Ka X-ray line of Pb atom. [Given Z for Pb = 82 and R¢ = 3.289 ¥ 1015 sec–1
[Hint : R¢ = Rc]. [Ans: 1.618 ¥ 1019 Hz]
Nanoscience and
Nanotechnology
22
Learning Objectives
pyrolysis, their properties, inorganic
After reading this chapter you will be able to nanotubes and biopolymers
LO 1 Differentiate between nanomaterials LO 3 Know nanoscales in 2D and 3D,
and bulk materials, and nanoscience nanoparticles and various methods of
and nanotechnology their synthesis including Ball milling,
gas condensation, sputtering, CVD,
LO 2 Understand quantum confinement,
CVC, Sol-gel and electrodeposition
nanowires and their synthesis through
techniques, properties of nanoparticles,
top down and bottom up approaches,
Bucky balls or fullerenes, their synthe-
VLS and VS methods, and catalyst
sis, properties and applications, QDs and
free growth mechanism, single walled
their fabrication and characterization
and multi-walled carbon nanotubes
and their fabrication methods, i.e., LO 4 Explain applications of nanotechnology
arc discharge method, laser ablation LO 5 Evaluate limitations and disadvantages
method, CVD technique, n-hexane of nanotechnology
Introduction
The prefix nano in nanotechnology means a billionth (1 ¥ 10–9 ∫ 1 nm). The typical dimension of
nanomaterials or nanostructures spans from subnanometer to several hundred nanometers. Generally
the dimesion/length may be less than 100 nm. Figure 22.1 shows how things scale and how small a
nanometer is. Depending on the number of dimension, the nanomaterials are classified as quantum dots
(0D: zero dimension), quantum wires (1D: one dimension) and quantum wells (2D: two dimensions).
Here the dimensionality, refers to the number of degrees of freedom in the electron momentum.*
In semiconductor industry, the efforts have been to reduce the size of the devices and the continued
decrease in device dimensions has followed the well-known Moore’s Law. The Moore’s law was predicted
in 1965, which states that the transistor size is decreased by a factor of 2 in every 18 months. In view of
this, it is said that the study of materials in the nanometer scale is partly driven by the ever shrinking of
devices in the semiconductor industry.
*
If the number of degrees of freedom is 3, then the material is said to be in bulk form. Since the electrons are free to move in three
dimensions, the extent of the confinement in bulk materials is zero. Similarly the extent of the confinement in quantum well is 1, in
quantum wire it is 2 and in quantum dot it is 3.
22.1 HOw NaNOmateriaLs are DiffereNet frOm BuLk materiaLs? LO1

The difference in the properties of nanostructures and that of micro or larger scale systems has its origin in
a very fundamental physical property of materials. In any material an atom on the surface is different from
an atom in the bulk. This difference is primarily because of the different potentials that the two atoms are
experiencing. Now consider as a simple case a sphere of radius r, the volume of this sphere is V = 4 p r 3 and
3
the surface area is given by S = 4pr2. So the ratio of the atoms on the surface to that in the bulk can roughly
be approximated as,
Nb r
=
Ns 3
where Nb is the number of atoms in the bulk and Ns is the number of atoms on the surface. Now this
approximation gives us an idea that as the size of the particle decreases the number of atoms on the surface is
a greater proportion of the total number of atoms. Hence, as the particle size becomes smaller and smaller the
surface atoms start dominating the properties of the whole material. This is what gives rise to the difference
in the properties of nanoparticles and their so varied applications.
At the nano-scale the quantum effects become important. For example, confinement of carriers, discrelization
of energy level and modification in band gap take place. The high surface area to volume ratio enhances
the ability of the structure to take part in chemicals reaction. Due to the absence of dislocations between
separate crystals the nonomaterials carry high tensile strength. The ability of manipulating the dimensions
of a structure allows one to tune the optical properties. Their novel optical properties enable them to absorb
or emit precise frequencies of light. Finally the simple difference is the size effect which plays the important
role in the applications of nanomaterials.
22.2 DiffereNce BetweeN NaNOscieNce aND NaNOtecHNOLOgy LO1

Nanoscience is the study of objects having size less than hundred nanometers at least in one dimension,
whereas nanotechnology is the engineering of these objects using different techniques. Nanotechnology is all
about the techniques and tools to come up with a nanoscale design or system, which exploit the properties at
molecular level to be more accurate and efficient. The difference between nanoscience and nanotechnology
is just like the difference between theory and practice. If we say that science has developed it does not
necessarily mean that the technology has developed. However, if we say that the technology has developed,
it means science might have definitely developed. Nanotechnology deals with the ability to develop and use
the technology to manipulate and observe at the nanoscale, whereas nanoscience is the study of phenomena
and objects at the nanoscale. We can say that nanoscience is just the study of nanostructures, whereas
nanotechnology is the application and implementation of understandings in different industries. We can take
an example of carbon nanotubes which are a subject of nanoscience and we will have to work on small scale
to create them. However, these carbon nanotubes will be of great use as a building material and the study
about its usage can be referred to as nanotechnology. Nanotechnology is applied in many areas including
information technology, automobile, textile, healthcare and agricultural industries.
22.3 QuaNtum cONfiNemeNt aND cLassificatiON Of LO2

NaNOstructures
Quantum confinement is the restricted motion of randomly moving electrons in specific energy levels when
the dimensions of a material approach the deBroglie wavelength of electron. The length scale corresponding to
Nanoscience and Nanotechnology 753
quantum confinement regime usually ranges from 1-25 nm. Quantum confinement leads to change in optical
and electrical properties. The change in properties occurs because of the energy levels becoming discrete and
because of restricted motion of electrons.
Based on the number of dimensions that are confined, nanostructures are classified as 3D (in which no dimen-
sions are confined); 2D (in which one dimension is confined, e.g. thin films); 1D (in which two dimensions are
confined; e.g. nano wires) and OD (in which all the dimensions are confined e.g. Quantum Dots, Nano crystals).
22.4 NaNoscale IN 1-D LO2

One-dimensional nanomaterials are nanofiber, such as nano rod, and nano tubular materials with fibre, such
as rod and tube, having length from 100 nm to tens of microns. Below we discuss in detail the nanowires,
carbon nanotubes, inorganic nanotubes and bio-polymers falling within this category.
22.4.1 Nanowires and their Different kinds

A nanowire is a nanostructure having the diameter of the order of a nanometer. A structure with the ratio of
the length to width greater than 1000 is also termed as nanowire. This ratio is known as aspect ratio. We can
say that the nanowires are the structures which have a thickness (diameter) constrained to 10 nm or less and
an unconstrained length. Since these structures have an unconstrained longitudinal size, these are referred to
as one-dimensional (1-D) materials. At these scales or if the diameter of the nanowire is of the order of the
de Broglie wavelength of the electron, then quantum effects become dominated. In view of this such wires
are also known as quantum wires. Nanowires have many interesting properties which are not seen in bulk
materials. The reason for this is that the electrons in nanowires are quantum confined laterally because of
which they occupy energy levels different from the traditional continuum of energy levels or bands found in
bulk materials. A quantum mechanical restraint on the number of electrons that can travel through the wire
at the nanometer scale leads to the discrete values of the electrical conductance, which are often referred to
as the quantum of conductance.
Nanowires can be superconducting, metallic, semiconducting and insulating in nature. For example, the
wires prepared using YBCO are superconducting in nature, whereas the nanowires made from SiO2 or TiO2
are insulating. The nanowires of Ni, Pt and Au behave as metals, while the silicon nanowires (SiNWs), InP
and GaN wires are semiconducting in nature. Other type of the nanowires is the molecular nanowires which
are composed of repeating molecular units of either organic (DNA) or inorganic (Mo6S9-xIx).
There exit many applications where nanowires may become important in electronic, optoelectronic and
nanoelectromechanical devices. These are used as additives in advanced composites. These are used as
field-emitters and as leads for biomolecular Nano sensors. The nanowires are also used in making metallic
interconnects in nanoscale quantum devices.
22.4.1.1 Synthesis of Nanowires
A nanowire is an extremely thin wire. Two very simple processes by which nanowires can be manufactured
are suspension and deposition. In the method of suspension, a suspended nanowire is held up by the ends
in an evacuated chamber. Then, in order to reduce its diameter, the wire is chemically etched or bombarded
with high speed atoms or molecules. Another method, called the deposition method, involves indenting the
surface of a wire in the center of a suspended span, raising the temperature, and then stretching the wire while
it is near its melting point. This way the deposited nanowire is fabricated on a surface consisting of some
non-conducting substance such as glass or plastic. This process is similar to that by which semiconductor
chips are grown. The only difference is that here we get a linear (one-dimensional) structure rather than a flat
(two-dimensional) or solid (three-dimensional) structure.
In addition to the above processes, other two basic approaches to synthesizing nanowires are called top down
technique and bottom up technique.
(i) Top-Down Approach
In top-down approach materials are reduced from two-dimensional thin films to a desired structure.
In order to fabricate nanowires, lithographic techniques (photo lithography and electron lithography)
and various etching techniques are used.
(ii) Bottom-Up Approach
Since nanowires are a result of anisotropic 1-D crystal growth on a nanometer scale, the key issue
related to the growth of nanowires is how to induce 1-D crystal growth in a controlled manner.
For accomplishing this many approaches have been employed, which include the metal catalyst
assisted vapour liquid solid (VLS) mechanism, vapour solid (VS) mechanism and template assisted
mechanism.
Si precursors
(vapour)
Au–Si alloy (liquid)
Si precursors Si crystal (solid)

Au particle (vapour)
Au–Si alloy
(liquid)
Si substrate
(i) Solid catalyst (ii) Alloying and melting (iii) Nucleation and growth
Figure 22.1
(a) Vapour Liquid Solid (VLS) Method

This is a common technique for creating a nanowire, which is used to grow single crystalline
nanowires of semiconducting materials. It requires a catalyst particle. A typical VLS process starts
with the dissolution of gaseous reactants into nano-sized liquid droplets of a catalyst metal, followed
by nucleation and growth of single crystalline rods and then wires. In this mechanism, the grown
species are supplied from vapour to metal catalyst that seat on the nanowire and remain on the top
of the nanowire during growth. The one-dimensional growth is induced and dictated by the liquid
droplets, whose sizes remain essentially unchanged during the entire process of wire growth. Here
each liquid droplet serves as a virtual template to strictly limit the lateral growth of an individual
wire. The growth is done at or above the eutectic point of the catalyst metal particle. A schematic
diagram of VLS growth is shown in Figure 22.1.
(b) Vapour Solid (VS) Method
Vapour solid (VS) mechnism is generally referred to as a process of one-dimensional material
growth from vapour phase precursors in the absence of the catalyst or obvious VLS evidence. In this
technique vapour is first generated by evapouration, chemical reduction or gaseous reactions. Then
vapour is subsequently transported and condensed on the substrate. This enables one to grow the
1-D nanostructures using the VS process, if one can control its nucleation and its subsequent growth
process. In consideration of thermodyanamics and kinetics, the VS growth of nanowires could be
possible via (i) a self catalytic VLS growth, (ii) an oxide assisted growth and (iii) Frank’s screw
dislocation mechnism.
(c) Catalyst Free Growth Mechnism
Through this mechanism the self assembles growth and the patterned growth of nanowires can be
achieved. In self assembled growth, the nanowires grow by self assembly without any lithographic
technique and without metal particle. In this mechnism, an oxide layer of SiO2 is deposited over
the substrate, usually Si. The SiO2 is amorphous in nature and contains large number of pin holes
which act as nucleation center for the growth of nanowires. The nanowires grow at an optimum
thickness of the SiO2. Then vapours of growing materials are supplied by various techniques such
as molecular beam epitaxy (MBE), metal organic chemical vapour deposition (MOCVD) and pulse
laser deposition (PLD). Each technique has its own advantages and some disadvantages. Since the
grown nanowires contain the same orientation as the substrate in this technqiue, the growth is called
epitaxial growth of nanowires. On the other hand, in Patterned growth, a pattern is made on the
oxide grown substrate by electron beam lithography. Then vapours are provided by choosing the
appropriate technique. Finally, the nanowires grow at the holes only (Figure 22.2). So by patterning
the substrate we can grow an ordered array of nanowires which can be used in many applications.
(a) (b)
(c) (d)
Figure 22.2
The nanowires can be characterized for their structural, transport and optical properties. The
characterization techniques used for the structural properties are scanning electron microscopy
(SEM), transmission electron microscopy (TEM), scanning tunneling probes (STPs), X-ray
analysis and Raman spectroscopy. Their transport properties are characterized with the help of I-V
measurement, temperature dependent resistance measurements and magnetoresistance. The optical
properties of these nanowires can be investigated through photoluminescence.
22.4.2 Carbon Nanotubes
Carbon nanotubes were discovered in 1991 by S. Iijima. These opened up a new era in materials science.
Carbon nanotubes, abbreviated as CNTs, are allotrope of carbon with a cylindrical nanostructure as thin as
a few nanometers yet be as long as hundreds of microns. Owing to their extraordinary thermal conductivity
and mechanical and electrical properties, carbon nanotubes find application as additives to various structural
materials in addition to their importance in nanotechnology electronics and optics. Nanotubes are members
of the fullerene structural family. The structure of a carbon nanotube is formed by a layer of carbon atoms
that are bonded together in a hexagonal (honeycomb) mesh. This one atom thick layer of carbon is called
graphene, which is wrapped in the shape of a cylinder and bonded together to form a carbon nanotube. The
combination of the rolling (chiral) angle and radius decides the nanotubes properties; for example, whether
the individual nanotube shell is a metal or semiconductor. Carbon nanotubes have a range of electric, thermal,
and structural properties that can change based on the physical design of the nanotube. Nanotubes can have
a single outer wall of carbon called single-walled nanotubes (SWNTs) or they can be made of multiple
walls (cylinders inside other cylinders of carbon) called multiple-walled nanotubes (MWNTs). Figure 22.3
shows how various types of carbon nanotubes can be formed from graphene. Similar to those of graphite, the
chemical bonding of nanotubes is composed entirely of sp2 bonds. These bonds which are stronger than the
sp3 bonds found in alkenes and diamond provide nanotubes with their unique strength.
Carbon nanotubes have found applications in field emitters / emission, conductive or reinforced plastics,
energy storage, molecular electronics with CNT based nonvolatile RAM and transistors. The CNT based
ceramics, fibers and fabrics have also attracted the people for their enormous usage.
Single wall carbon nanotube Multi walled carbon nanotube

Figure 22.3
22.4.2.1 Kinds of Carbon Nanotubes

A broad category of carbon nanotubes comprises single-walled carbon nanotubes and multi-walled carbon
nanotubes.
(i) Single-walled Carbon Nanotubes (SWCNTs)
Single-walled carbon nanotubes can be formed in three different designs, namely, Armchair, Chiral
and Zigzag (Figure 22.4). The design depends on the way the graphene is wrapped into a cylinder.
For example, rolling a sheet of paper from its corner will assign one design, and rolling the paper
from its edge will attain a different design. The structure of a SWCNT is represented by a pair of
indices (n, m) called the chiral vector. The structural design has a direct impact on the electrical
properties of the nanotube. When n − m is a multiple of 3, then the nanotube is described as metallic,
i.e., highly conducting, otherwise the nanotube is a said to behave as a semiconductor. The Armchair
design is always metallic while other designs can make the nanotube a semiconductor.
(ii) Multi-walled Carbon Nanotubes (MWCNTs)

There are two structural models of multi-walled nanotubes, namely, Russian doll model and Parch-
ment model. In the Russian doll model, a carbon nanotube contains another nanotube inside it such
that the inner nanotube has a smaller diameter than the outer nanotube. In the Parchment model, a
single graphene sheet is rolled around itself multiple times, resembling a rolled up scroll of paper.
Multi-walled carbon nanotubes have similar properties to single-walled nanotubes. However, the outer
walls of these nanotubes can protect the inner carbon nanotubes from chemical interactions with out-
side materials. These nanotubes also have a higher tensile strength than the single-walled nanotubes.
22.4.2.2 Fabrication of Carbon Nanotubes
There are several procedures that have been developed for fabricating carbon nanotube structures.
(a) Arc Discharge Method
Arc discharge method is used for the synthesis of carbon nanotubes where a DC arc voltage is
applied across two graphite electrodes (anode and cathode) immersed in an inert gas (such as He)
filling a chamber. The electrodes are vapourized by the passage of a DC current of about 100 A
through the two high-purity graphite separated by about 1–2 mm in 400 mbar of He atmosphere. In
this process a carbon rod is built up at the cathode after arc discharging for a period of time. This
method can mostly produce multi-walled nanotubes but can also produce single-walled nanotubes
with the addition of metal catalyst such as Fe, Ni, Co, Y or Mo, on either the anode or the cathode.
Various parameters such as the metal concentration, inert gas pressure, type of gas, plasma arc,
temperature, the current and system geometry are the deciding factors for the quantity and quality
such as lengths, diameters and purity of the nanotubes. When pure graphite rods are used, fullerenes
are deposited as soot inside the chamber and multi-walled carbon nanotubes are deposited on the
cathode. However, single-walled carbon nanotubes are generated in the form of soot when a graphite
anode containing a metal catalyst (Fe or Co) is used with a pure graphite cathode. The advantage of
this method is that it produces a large quantity of nanotubes. But the main disadvantage is that there
is relatively a little control over the alignment (i.e., chirality) of the produced nanotubes, which is
critical to their characterization and role.
(b) Laser Ablation Method
Laser ablation process was developed by Richard Smalley and co-workers at Rice University,
who were primarily involved in blasting metals with a laser to produce various metal molecules.
However, they replaced the metals with graphite to create multi-walled carbon nanotubes when
they heard of the existence of nanotubes. In laser ablation, a pulsed laser is used to vapourize a
graphite target (kept in a quartz tube) in a high-temperature reactor, where a flow of argon gas
is maintained throughout the reaction. Here the carbon vapourizes, which is carried away by the
argon gas and is condensed downstream on the cooler walls of the quartz. This condensation leads
to the formation of single-walled carbon nanotubes and metallic particles. Therefore, purification
methods are employed to this mixture. The key to the proper formation of the condensed nanotubes
is that the location where the carbon atoms begin to condense should be set up as a curved sheet of
graphene with a catalyst metallic atom nearby. As carbon atoms begin to attach and form rings, the
metallic atom, if it has the proper electronegativity properties, will preserve the open edge of the
tube and prevent it from drawing to a close. This phenomenon is termed as the scooter effect. The
laser ablation method yields around 70% and produces primarily single-walled carbon nanotubes
with a controllable diameter, which can be determined by the reaction temperature. This technique
includes a relatively low metallic impurities, since the metallic atoms involved tend to evaporate
from the end of the tube once it is closed. However, it is more expensive than either arc discharge or
chemical vapour deposition techniques. Moreover, the nanotubes produced from this method are not
necessarily uniformly straight, but instead do contain some branching.
(c) Chemical Vapour Deposition Technique
The chemical vapour deposition (CVD) technique allows carbon nanotubes to grow on a variety of
materials, which makes it more viable to integrate into already existent processes for synthesizing
electronics. This process involves the chemical breakdown of a hydrocarbon on a substrate, which is
due to the fact that a main way to grow carbon nanotubes is by exciting carbon atoms that are in contact
with metallic catalyst particles. The CVD method extends this idea by embedding these metallic
particles (for example, iron) in properly aligned holes in a substrate (say, silicon). Essentially, tubes
are drilled into silicon and implanted with iron nanoparticles at the bottom. After that a hydrocarbon
such as acetylene is heated and decomposed onto the substrate. The carbon comes into contact with
the metal particles embedded in the holes and start to form nanotubes, which are templated from
the shape of the tunnel. The advantages of this method are that the yield is very high and the size
of the growth area is theoretically arbitrary. Moreover, the alignment of the nanotubes is consistent,
which is crucial for creating particular types of nanotubes, e.g. semiconductor or metallic. However,
the main disadvantage is that the large sized areas (several millimeters) tend to crack, shrink, and
otherwise twist. Hence, the substrates need to be dried very thoroughly to prevent against this.
(d) n-hexane Pyrolysis
Researchers have developed a method to synthesize large, long single-walled nanotube bundles in a
vertical furnace by pyrolyzing hexane molecules. These n-hexane molecules are mixed with certain
other chemicals, which have been shown independently to help with growth of nanotubes. These
are burned (pyrolyzed) at a very high temperature in a flow of hydrogen and other optional gases.
The use of a different hydrocarbon or using a different gas has been shown to prevent the formation
of long nanotubes. The primary advantage of this method is that it produces macroscopic nanotube
bundles (micro tubes), whose their diameters are typically larger than that of human hair, and their
length is several centimeters. However, the disadvantage is that the alignment is not as produced
from other methods, making it viable for creating micro cables, but not nanotubes with precise
electrical properties. Moreover, the elasticity of these nanotube bundles is not found to be as great
as hoped, i.e., they are more brittle.
22.4.2.3 Properties of Carbon Nanotubes
The carbon nanotubes are known for their strength, electrical properties and thermal properties. These
properties are discussed below.
(a) Strength
Carbon nanotubes have a higher tensile strength than steels, which comes from the sp² bonds between
the individual carbon atoms. This bond is even stronger than the sp³ bond found in diamond. Under
high pressure, individual nanotubes can bond together and exchange some sp² bonds with sp³ bonds.
This way the possibility of producing long nanotube wires is enhanced. Carbon nanotubes are not
only strong but they are also elastic. However, their elasticity does have a limit, it is possible to
permanently deform to shape of a nanotube under very strong forces. The strength of a nanotube
is weakened by defects in its structure, occurring from atomic vacancies or a rearrangement of the
carbon bonds. This in turn causes the tensile strength of the entire nanotube to weaken.
(b) Electrical Properties

The structure of a carbon nanotube determines how conductive the tube is. A carbon nanotube is
highly conductive under the situation when the structure of atoms in it minimizes the collisions
between the conduction electrons and the atoms. The strong bonds between carbon atoms also allow
carbon nanotubes to withstand higher electric currents than copper. The electron transport occurs
only along the axis of the nanotube. Single-walled nanotubes can route electrical signals at the
speeds up to 10 GHz when used as interconnects on semiconducting devices.
(c) Thermal Properties
The strength of the atomic bonds in carbon nanotubes allows them to withstand high temperatures,
because of which carbon nanotubes have been shown to be very good thermal conductors. Compared
with copper wires which are commonly used as thermal conductors, the carbon nanotubes can
transmit over 15 times the amount of Watts/m/K. However, the thermal conductivity of carbon
nanotubes is dependent on the temperature of the tubes and the outside environment.
22.4.3 Inorganic Nanotubes
An inorganic nanotube is a cylindrical molecule that is often composed of metal oxides or group III-Nitrides.
These nanotubes are morphologically similar to carbon nanotubes, but are non-carbon nanotubes. The inorganic
nanotubes have been found naturally in some mineral deposits the way carbon nanostructure are found naturally.
Minerals such as white asbestos and imogolite have been shown to have a tubular structure. The inorganic nano-
tubes have been synthesized based on molybdenum disulphide (MoS2) and tungsten disulphide (WS2).
Inorganic nanotubes are the nanostructures which are distinct from pure inorganic nano wires and carbon nanotubes
in two ways. First, the resultant physical properties and electronic structure show combined characteristics of both
one- and two-dimensional materials. Second, these hollow nanotubes can serve as nanoscale containers or pipes
to deliver fluids and molecular species. These are excellent building blocks for the construction of large scale nano
fluidic systems. Inorganic nanotubes show easy synthetic access and high crystallinity, needle like morphology,
good uniformity and dispersion, and good adhesion to a number of polymers. So these are promising candidates as
fillers for polymer composites with enhanced thermal, mechanical and electrical properties. Inorganic nanotubes
are heavier than carbon nanotubes and are not as strong under tensile stress, but these are particularly strong under
compression, leading to potential applications in impact resistant applications such as bulletproof vests.
22.4.4 Biopolymers
Biopolymers are polymers which are produced by living organisms. It means these are the polymers that are
biodegradable. In other words, we can say that these are polymeric biomolecules. The biopolymers represent
the most abundant organic compounds in the biosphere and constitute the largest fraction of cells. Their main
classes are distinguished according to their chemical structures. For example, there are four main types of bio-
polymer based on sugar, starch, cellulose and synthetic materials. The input materials to produce these polymers
may be either synthetic or renewable, i.e., based on agricultural plant or animal products. Two main strategies
may be followed in synthesizing a polymer. One is to build up the polymer structure from a monomer. This
is done by a process of chemical polymerization. The alternative is to take a naturally occurring polymer and
chemically modify it to give it the desired properties. However, the biodegradability of the polymer may be
adversely affected by chemical modification. Therefore, it is often necessary to seek a compromise between the
desired material properties and biodegradability.
Some biopolymers such as polylactic acid (PLA), naturally occurring zein, and poly-3-hydroxybutyrate can
be used as plastics, which replace the need for polystyrene or polyethylene based plastics. Some plastics are
now referred to as being degradable, oxydegradable or UV degradable, which means they break down when
exposed to air or light. However, these plastics are still largely oil based.
For their use in the packaging industry as food trays, blown starch pellets for shipping fragile goods, thin
films for wrapping, etc., biopolymers are produced from biomass which comes from crops such as sugar beet,
wheat or potatoes. The conversion of these takes place in the following manner. For example, sugar beet is
converted to glyconic acid which finally gives polyglyconic acid. The fermentation of starch is done in order
to get lactic acid which is converted to polylactic acid. On the other hand, the fermentation of biomass leads
to the bioethanol which gives ethene and the ethene is converted into polyethylene.

Thin films with thickness ranging from atomic monolayer to a few atomic layers are 2D nanomaterials. They
are characterized by a high aspect ratio and their small size perpendicular to the film gives rise to quantum
confinement of electrons.
Synthesis of thin films is carried out by either physical or chemical techniques. The chemical techniques
employed are chemical vapour deposition (CVD), Sol-Gel technique and electrochemical plating. Chemical
techniques are based on decomposition of molecular precursors and subsequent deposition of the desired
material on the target.
Physical techniques are based on the vaporization of material to be deposited and its condensation n the
substrate surface as thin film. Based on the method used for vaporization, physical techniques are classified
as Molecular Beam Epitaxi (MBE), sputtering, and Pulsed Laser Ablation and Deposition (PLD).

If there is no restriction of carrier confinement or carriers are free to move in all 3 directions, the material
is called bulk material (also termed as 3-D). In this case the property or the behavior of nanomaterials stand
the same as that of the materials. It means the energy levels are continuous and density of state varies as E3/2,
where E is the energy of the state. Ingot GaAs, Si, Ge etc. fall within this category of 3-D nanomaterials.
22.6.1 Nanoparticles
Figure 22.4 shows how things scale and how small a nanometer actually is. Although it may seem that such
structures have come into being in the very recent past, this is not true. Humans have been known to take
advantage of the peculiar properties of nanoparticles as early as the 4th century A.D. Roman glassmakers
were fabricating glasses containing nano sized metals. The great varieties of beautiful colors of the windows
of medieval cathedrals are due to the presence of metal nanoparticles in the glass.
22.6.1.1 Synthesis of Nanoparticles
Nanoparticles can be synthesized by means of various techniques, some of which are discussed below.
22.6.1.1.1 Mechanical Method Ball milling is the best example for mechanical method. A ball mill
(Figure 22.5), a type of grinder, is a cylindrical device used in grinding (or mixing) materials like ores,
chemicals, ceramic raw materials and paints. Ball mills rotate around a horizontal axis, partially filled with
the material to be ground plus the grinding medium. Different materials are used as media, including ceramic
balls, flint pebbles and stainless steel balls.
Figure 22.4
An internal cascading effect reduces the material to a fine powder.

Industrial ball mills can operate continuously fed at one end and
discharged at the other end. Large to medium-sized ball mills are
mechanically rotated on their axis, but small ones normally consist of
a cylindrical capped container that sits on two–drive shafts (pulleys
and belts are used to transmit rotary motion).
A rock tumbler functions on the same principle. Ball mills are also
used in pyrotechnics and the manufacture of black powder, but cannot
be used in the preparation of some pyrotechnic mixtures such as flash
powder because of their sensitivity to impact. High-quality ball mills
are potentially expensive and can grind mixture particles to as small
as 0.0001 mm, enormously increasing surface area and reaction rates.
Figure 22.5
The grinding works on principle of critical speed. The critical speed
can be understood as that speed after which the steel balls (which are responsible for the grinding of particles)
start rotating along the direction of the cylindrical device; thus cause no further grinding.
22.6.1.1.2 Gas Condensation Nanocrystalline metals and alloys were first fabricated using the gas
condensation technique. In this technique, a metallic or inorganic material is vaporized using thermal
evaporation sources such as a refractory crucibles or electron beam evaporation devices, under a pressure
of 1–50 mbar. In gas evaporation, a high residual gas pressure causes the formation of ultra fine particles
(100 nm) by gas phase collision. The ultra fine particles are formed by collision of evaporated atoms with
residual gas molecules. The sources used for vaporization may be resistive heating, high energy electron
beams, low energy electron beam and inducting heating. Clusters form in the vicinity of the source by
homogenous nucleation in the gas phase that grows by incorporation by atoms in the gas phase. It comprises
a ultra high vacuum system fitted evaporation source, a cluster collection device of liquid nitrogen filled cold
finger scrapper assembly and a compaction device. During heating, atoms condense in the super saturation
zone close to the heating device. The nanoparticles are removed by scrapper in the form of a metallic plate.
Evaporation is to be done from W, Ta or Mo refractory metal crucibles. In case the metals react with crucibles,
we can use the electron beam evaporation technique. The method is extremely slow. A few limitations of the
method are source precursor incompatibility, temperature ranges and dissimilar evaporation rates in an alloy.
22.6.1.1.3 Sputtering Alternative sources have been developed over the years. For instance, Fe is evaporated
into an inert gas atmosphere. Through collision with the atoms the evaporated Fe atoms loose kinetic energy
and condense in the form of small crystallite crystals that accumulate as a loose powder. Sputtering or laser
evaporation may be used instead of thermal evaporation. Sputtering is a nonthermal process in which surface
atoms are physically ejected from the surface by momentum transfer from an energetic bombarding species
of atomic/molecular size. Typical sputtering uses a glow discharge or ion beam. Interaction events which
occur at and near the target surface during the sputtering process in magnetron sputtering has advantage
over diode and triode sputtering. In magnetron sputtering, most of the plasma is confined to the near target
region. Other alternate energy sources which have been successfully used to produce clusters or ultra fine
particles are sputtering electron beam heating and plasma methods. Sputtering has been used in low pressure
environment to produce a variety of clusters including Ag, Fe and Si.
22.6.1.1.4 Vacuum Deposition and Vaporization Before proceeding to the other methods, it is
important to understand the terms vacuum deposition and vaporization or vacuum evaporation. In vacuum
deposition process, elements, alloys or compounds are vaporized and deposited in a vacuum. The vaporization
source is the one that vaporizes materials by thermal processes. The process is carried out at pressure of less
than 0.1 Pa (1 m Torr) and in vacuum levels of 10 to 0.1 MPa. The substrate temperature ranges from room
temperature to 500°C. The saturation or equilibrium vapour pressure of a material is defined as the vapour
pressure of the material in equilibrium with the solid or liquid surface. For vacuum deposition, a reasonable
deposition rate can be obtained if the vaporization rate is fairly high. Vapour phase nucleation can occur
in dense vapour cloud by multi body collisions, the atoms are passed through a gas to provide necessary
collision and cooling for nucleation. These particles are in the range of 1 to 100 nm and are called ultra
fine particles or clusters. The advantages associated with vacuum deposition process are high deposition
rates and economy. However, the deposition of many compounds is difficult. Nanoparticles produced from a
supersaturated vapour are usually longer than the cluster.
22.6.1.1.5 Chemical Vapour Deposition (CVD) and Chemical Vapour Condensation (CVC)
Chemical Vapour Deposition (CVD) is a well–known process in which a solid is deposited on a heated
surface via a chemical reaction from the vapour or gas phase. CVD reaction requires activation energy to
proceed. This energy can be provided by several methods. In thermal CVD, the reaction is activated by a
high temperature above 900°C. A typical apparatus comprises gas supply system, deposition chamber and
an exhaust system. In plasma CVD, the reaction is activated by plasma at temperatures between 300°C and
700°C. In laser CVD, pyrolysis occurs when laser thermal energy heats an absorbing substrate. In photo laser
CVD, the chemical reaction is induced by ultra violet radiation which has sufficient photon energy, to break
the chemical bond in the reactant molecules. In this process, the reaction is photon activated and deposition
occurs at room temperature. Nano composite powders have been prepared by CVD. SiC/Si3N composite
powder was prepared using SiH4, CH4, WF6 and H2 as a source of gas at 1400°C.
Another process called Chemical Vapour Condensation (CVC) was developed in Germany in 1994. It
involves pyrolysis of vapours of metal organic precursors in a reduced pressure atmosphere. Particles of
ZrO2, Y2O3 and nano whiskers have been produced by CVC method. A metal-organic precursor is introduced
in the hot zone of the reactor using mass flow controller. For instance, hexamethyldisilazane (CH3)3 Si NHSi
(CH3)3 was used to produce SiCxNyOz powder by CVC technique. The reactor allows synthesis of mixtures
of nanoparticles of two phases or doped nanoparticles by supplying two precursors at the front end of reactor
and coated nanoparticles, n-ZrO2, coated with n-Al2O3 by supplying a second precursor in a second stage of
reactor. The process yields quantities in excess of 20 g/hr. The yield can be further improved by enlarging the
diameter of hot wall reactor and mass of fluid through the reactor.
22.6.1.1.6 Sol-gel Techniques In addition to techniques mentioned above, the sol-gel processing
techniques have also been extensively used. Colloidal particles are much larger than normal molecules or
nanoparticles. However, upon mixing with liquid colloids appear bulky whereas the nanosized molecules
always look clear. It involves the evolution of networks through the formation of colloidal suspension (sol)
and gelatin to form a network in continuous liquid phase (gel). The precursor for synthesizing these colloids
consists of ions of metal alkoxides and aloxysilanes. The most widely used are tetramethoxysilane (TMOS),
and tetraethoxysilanes (TEOS) which form silica gels. Alkoxides are immiscible in water. They are organo-
metallic precursors for silica, aluminium, titanium, zirconium and many others. Mutual solvent alcohol is
used. The sol-gel process involves initially a homogeneous solution of one or more selected alkoxides. These
are organic precursors for silica, alumina, titania, zirconia, among others. A catalyst is used to start reaction
and control pH. Sol-gel formation occurs in four stages: (i) Hydrolysis (ii) Condensation (iii) Growth of
particles, and (iv) Agglomeration of particles.
(i) Hydrolysis
During hydrolysis, addition of water results in the replacement of [OR] group with [OH-] group.
Hydrolysis occurs by attack of oxygen on silicon atoms in silica gel. Hydrolysis can be accelerated
by adding a catalyst such as HCI and NH3. Hydrolysis continues until all alkoxy groups are replaced
by hydroxyl groups. Subsequent condensation involving silanol group (Si-OH) produced siloxane
bonds (Si-O-Si) and alcohol and water. Hydrolysis occurs by attack of oxygen contained in the
water on the silicon atom.
(ii) Condensation
Polymerization to form siloxane bond occurs by either a water producing or alcohol producing
condensation reaction. The end result of condensation products is the formation of monomer, diamer,
cyclic tetramer, and high order rings. The rate of hydrolysis is affected by pH, reagent concentration
and H2O/Si molar ratio (in case of silica gels). Also ageing and drying are important. By control of
these factors, it is possible to vary the structure and properties of sol-gel derived inorganic networks.
(iii) Growth and Agglomeration
As the number of siloxane bonds increases, the molecules aggregate in the solution, where they form
a network, a gel is formed upon drying. The water and alcohol are driven off and the network shrinks.
At values of pH of greater than 7, and H2O/Si value ranging from 7 to 5, spherical nanoparticles are
formed. Polymerization to form siloxane bonds by either an alcohol producing or water producing
condensate occurs.
2 HOSi (OR)3 Æ (OR)3 Si O Si (OR)3 + H2O
or
2 HOSi (OR)3 Æ (OR)2OH Si O Si (OR)3 + H2O
Above pH of 7, silica is more soluble and silica particles grow in size. Growth stops when the difference in
solubility between the smallest and largest particles becomes indistinguishable. Larger particles are formed
at higher temperatures. Zirconium and yttrium gels can be similarly produced.
Despite improvements in both chemical and physical methods of synthesis, there remain some problems and
limitations. Laser vaporization technique has offered several advantages over other heating techniques. A
high energy pulsed laser with an intensity flux of 106 – 107 W/cm2 is forced on target material. The plasma
causes high vaporization and high temperature (10,000°C). Typical yields are 1014 – 1015 atoms from the
surface area of 0.01 cm2 in a 10–8 s pulse. Thus a high density of vapour is produced in a very short time
(10–8 s), which is useful for direct deposition of particles.
22.6.1.1.7 Electrodeposition Nanostructured materials can also be produced by electrodeposition. These
films are mechanically strong, uniform and strong. Substantial progress has been made in nanostructured
coatings applied either by DVD or CVD. Many other non-conventional processes such as hypersonic plasma
particle deposition (HPPD) have been used to synthesize and deposit nanoparticles. The significant potential
of nanomaterial synthesis and their applications is virtually unexplored. They offer numerous challenges to
overcome. Understanding more of synthesis would help designing better materials. It has been shown that
certain properties of nanostructured deposits such as hardness, wear resistance and electrical resistivity are
strongly affected by grain size. A combination of increased hardness and wear resistance results in a superior
coating performance.
22.6.1.2 Properties of Nanoparticles
Nanoparticles are generally considered to be a number of atoms or molecules bonded together such that the
dimension of the bonded entity is of the order of 100 nm. And since 1 nm is 10 Å, hence the critical limit
for nanoparticles size is 1000 Å. For the sake of completeness to this size based criterion for nanoparticles,
we may say that particles that are large enough so that not almost all the atoms are at the surface are still
nanoparticles. But it is quite evident that this size based scheme is arbitrary. For example, the heme molecule,
FeC34H32O4N4 that is found in hamoglobin contains around 75 atoms. Thus, a more convincing definition of
nanoparticles would be that they have sizes smaller than the critical lengths for certain physical phenomenon.
This critical length can characterize processes like electrical conductivity or excitonic processes. So, one
definition of nanoparticles of metals can be given by their scattering length, which is the distance that an
electron moves between two successive collisions with the vibrating atoms or impurities in the material. It is
below these critical lengths that the materials begin to demonstrate new physical or chemical phenomenon
that is not observed in bulk.
Now that we know what nanoparticles are let us study how the physical and chemical properties of materials
change when we enter the nanoparticles paradigm.
(i) Optical Properties
The optical properties of nanoparticles are markedly different from those of bulk. However, the
changes that are observed are quite different for different materials. We would talk about the metals
first and then go on to talk about the semiconducting materials.
(ii) Metals: Surface Plasmons
In the case of metals as the size of the particles decreases we start observing oscillations of electron
gas on the surface of nanoparticles. These oscillations are called surface plasmons. So, if the
nanoparticles are exposed to an electromagnetic wave (light) having a wavelength comparable to or
greater than the size of the nanoparticles and the light has a frequency close to that of the surface
plasmon then the surface plasmon would absorb energy. Thus nanoparticles start exhibiting different
colors as their size changes and the frequency of the surface plasmon changes with it. This kind
of a phenomenon is not observed strongly in bulk metals. This frequency of the surface plasmon
absorption is a function of the dielectric constant of the material, size of the particles and also the
specific geometrical shape that the particle has. This phenomenon of surface plasmon resonance
and subsequent absorption was used to obtain different colors of the stained glasses used in the
medieval cathedrals. Surface plasmons have been used to enhance the surface sensitivity of several
spectroscopic measurements including fluorescence, Raman scattering, and second harmonic

generation. However, in their simplest form, SPR reflectivity measurements can be used to detect
DNA or proteins by the changes in the local index of refraction upon adsorption of the target
molecule to the metal surface. If the surface is patterned with different biopolymers, the technique
is called Surface Plasmon Resonance Imaging (SPRI).
(iii) Semiconductors
In the case of semiconducting nanoparticles the properties change in a different fashion. One of the
important properties that changes as the size of the nanoparticles changes is the absorption spectrum
of the material. In any semiconducting material there is a band gap between the conduction and
valence bands in the solid. This band gap is specific to materials like bulk Si has a band gap of 1.1
eV. The band gap in any material affects its electrical as well as optical properties. Taking an optical
perspective whenever light of any particular wavelength is passed through a material sample there
is some degree of absorption by the material. The strength of absorption depends on the material
and the wavelength passed. A plot of the absorption strength against the wavelength is what is
referred to as the absorption spectrum. Figure 22.6 shows the absorption spectra of a direct as well
as indirect band gap semiconductor. For a given material in its bulk state the absorption spectrum is
unique. But when the material is in the form of nanoparticles then the absorption spectrum changes
and undergoes a blue shift. By blue shift we mean that the whole spectrum shifts towards a smaller
wavelength.
Absorption Absorption
Transparent
region
Onset of direct Onset of indirect

photon photon transition
transition Onset of direct
transition
Crystal With Direct Gap Crystal With Indirect Gap

Figure 22.6
22.6.2 carbon allotropes, Bucky balls or Fullerenes

Fullerene is a molecule of carbon which could be in the form of a hollow sphere, ellipsoid, tube and other
shapes. In 1985, a new allotrope of carbon was discovered in which the atoms are arranged in closed shells.
The new form was found to have the structure of truncated icosahedrons and was named Buckminsterfullerene
after the architect Buckminster Fuller who designed geodesic domes in the 1960’s. Previously six crystalline
forms of the element carbon were known, namely two kinds of graphite, two kinds of diamond, chaoite
and carbon (VI). Fullerenes are similar in structure to graphite, which are composed of stacked graphene
sheets of linked hexagonal rings; they may also contain pentagonal or sometimes heptagonal rings. Spherical
fullerenes are also known as Bucky balls, as they resemble the balls used in football, whereas the cylindrical
ones are called carbon nanotubes or Bucky tubes. Figure 22.7 provides a graphical representation of the
60-carbon atoms called C60 fullerene, whereas Figure 22.8 illustrates the arrangement of 540-carbon atoms
called icosahedral fullerene C540 which is another member of the family of fullerenes.
An infinite number of spherical fullerenes are believed to be able to exist, the known forms of which include
C-60, C-70, C-76, C-84, C-240 and C-540. All fullerenes consist of 12 pentagonal faces and a varying number
of hexagonal faces. In general, for a fullerene C-n, there will be 12 pentagonal faces and half of n minus 10,
i.e., n/2-10, hexagonal faces. It means the C-60 fullerene, which is called Bucky ball, has 12 pentagonal
faces (rings) and 20 hexagonal faces (rings), forming a spheroid shape with 60 vertices for 60 carbons. The
pentagonal rings sit at the vertices of an icosahedron such that no two pentagonal rings are next to each other.
The average C-C bond distance measured using nuclear magnetic resonance (NMR) is found to be 1.44 A°.
A diameter of 7.09 A° is calculated for the C-60 based on the fact that the C-C distance is equal to 1.40 A°
for the hexagon bonds and 1.46 A° for the pentagonal bonds length.
22.6.2.1 Synthesis of Fullerenes
The three main methods to synthesize single-walled carbon molecules, either Bucky balls or nanotubes are i)
electric arc discharge method, ii) laser ablation method and iii) chemical vapour deposition technique. These
have already been discussed in detail. Now these methods are discussed in view of the synthesis of fullerenes.
(i) Electric Arc Discharge Method
In this method, arcs of alternating or direct current are passed through the graphite electrodes kept
in an atmosphere of helium gas at approximately 200 Torr pressure (Figure 22.9). The graphite is
evaporated that takes the form of soot. This is dissolved in a nonpolar solvent. The solvent is dried
away and the C-60 and C-70 fullerenes are separated from the residue. This method yields up to 70%
of C-60 and 15% of C-70 at the optimal current, He pressure and flow rate.
(ii) Laser Ablation Method
Laser ablation method is one of the three methods of laboratory and industrial synthesis of Bucky
balls in addition to single-walled and multi-walled nanotubes. Laser vaporization is also used
for fullerene production. In a typical apparatus, a pulsed Nd:YAG laser is used as a laser source
operating at 532 nm and 250 mJ of power and the graphite target is kept in a furnace at 1200 °C.
(iii) Chemical Vapour Deposition Technique
This technique is based on the thermal cracking of a carbon containing gas (e.g. a hydrocarbon or
carbon-monoxide) in the presence of a catalyst. Hydrogen gas or an inert gas like Ar are used as vector
gas. Sometimes metallocenes such as ferrocene, nickelocene or cobaltocenes are used whose cracking
generates both the nanometric metallic catalysts as well as carbon for the formation of nanotubes.
Pressure control
system
Arc Chamber
Cathode
Initial Arc
Anode
Linear Drive
Figure 22.9
22.6.2.2 Properties of Fullerenes
The Bucky ball becomes more than twice as hard as its cousin, diamond, when compressed to 70 percent of
its original size. These can withstand slamming into a stainless steel plate at 15,000 mph, merely bouncing
back unharmed. This shows their high speed collisions property. The Bucky ball is the only known carbon
allotrope that can be dissolved at room temperature and for which aromatics are the best solvents. Larger
fullerenes (C-72) with trapped lanthanides have been found to have higher solubility.
22.6.2.3 Potential Applications of Fullerenes
Basic fullerenes and their functionalized derivatives have been suggested to have a large number of
applications. Potential applications include organic hydrogen gas storage, sensors, polymer electronics,
photovoltaic, as a molecular wire, as a precursor to diamond antioxidants, biopharmaceuticals, antibacterial,
HIV inhibition, catalysts, water purification, MRI agents, optical devices, scanning tunneling microscopy,
and atomic force microscopy. Fullerenes are being extensively investigated as carrier species for medical
radionuclides in cancer therapy. On the other hand, Bucky papers are used in fire resistance and in television
screens since these may be more efficient than CRT and LCD displays.
Pristine Fullerenes with no functional groups can have a positive effect and act as antioxidents. However,
functionalized fullerenes or fulleres dissolved in organic solvents are hazardous to the environments.
22.6.3 Quantum Dots (QDs)

Quantum dots (QDs) are semiconductor nanocrystals that have a collection of electrons varying from a
single electron to several thousands of electrons. Quantum dots are termed as zero- dimensional structured
materials, as the carriers, i.e., the electrons, are spatially confined in all the three directions or the motion of
the electrons is restricted in all the three dimensions. The energy levels of quantum dots are discrete and just
as an atom their energy levels are quantized and isolated. In view of this, the quantum dots are sometimes
referred to as artificial atoms. The density of state of quantum dots is a delta function of energy E. For this
reason, many electrons can be accommodated in a single energy level. This property makes the quantum dots
as an active component of high power and low threshold semiconductor LASERs and highly efficient photo
detectors in addition to their application to solar cells. A qualitative description of the response of quantum
dots of many shapes can be made based on the cuboid quantum dot, which is more often designated as the
quantum box and is a special case of zero-dimensional structures of other shapes such as spherical quantum
dots. The energy gap between the energy levels is found to be larger for the smaller quantum dots and hence,
the emitted light with higher frequency is observed.
22.6.3.1 Fabrication of QDs and Their Characterization
There are several techniques for the fabrication of quantum dots, such as molecular beam epitaxy (MBE),
metal organic chemical vapour deposition (MOCVD), pulsed laser deposition (PLD), etc. As in the case
of other nanostructures, the basic two approaches (top down approach and bottom up approach) are used
to fabricate the QDS. The top down approach involves lithography and etching in quantum well structure,
whereas the bottom up approach is related to the self-assembled growth of QDs.
The first step of the lithographic procedure used in top down approach is to place a radiation-sensitive resist
on the surface of the sample substrate. Then the sample is irradiated by an electron beam in the region where
the nanostructure is required to be formed. This can be done by using a radiation mask that contains the
nanostructure pattern or a scanning electron beam that strikes the surface only in the desired region. The next
step is the application of the developer to remove the irradiated portions of the resist. After that an etching
mask is inserted into the hole in the resist. Subsequently the remaining parts of the resist are lifted off. The
areas of the quantum structure not covered by the etching mask are then chemically etched away in order to
produce the quantum structure. Finally, the etching mask is removed.
Recent techniques for fabrication of quantum dots in bottom up approach involve strain induced self-assembly.
The term self-assembly represents a process where a strained 2-D system reduces its energy by changing into
a 3-D morphology. The InxGa1-xAs/GaAs system that offers a large lattice mismatch (7.2% between InAs
and GaAs) is employed as the material combination most commonly used for this technique. Through this
technique self-assembled InAs quantum dots can be grown on GaAs and the size, separation and height of the
quantum dots can be controlled by the deposition parameters. As was the case with quantum nanowires, the
random distribution of the quantum dots is, however, one of the drawbacks of this technique.
There are several techniques to characterize the quantum dots. These are atomic force microscope (AFM),
scanning tunneling microscope (STM), transmission electron microscope (TEM), photoluminescence (PL)
and in situ reflection high energy electron diffraction (RHEED).
22.7 appLicatiONs Of NaNOtecHNOLOgy LO4

Nanotechnology has been identified as the technology for the future and holds the promise to be the basic
tenet of the next industrial revolution. Although there are still existing applications of nanotechnology in the
world but the future holds much more promise. Some of the existing products that derive their functionality
from nanotechnology are discussed below.
22.7.1 self–cleaning Glass
Nanoparticles are coated on the glass to make it photocatalytic and hydrophilic. The photocatalytic effect
implies that when UV radiation from the light hits the glass the nanoparticles become energized and begin to
break down the organic particles on the glass surface and due to its hydrophilic nature the glass attracts water
particles, which then clean it.
22.7.2 clothing
Scientists are using nanoparticles to enhance tour clothing. By coating the fabrics with a thin layer of Zinc
Oxide nanoparticles clothes offering better protection from UV radiation can be created. Also, clothes can
have nanowhiskers that can make them repel water and other materials thus making them stain resistant.
Silver nanoparticles have been demonstrated to have an antibacterial effect on the clothes that they were
coated on to.
22.7.3 scratch Resistant coating

Materials like glass are being coated with thin films of hard transparent material to make the glass scratch
resistant. Anti–fog glasses with transparent nanostructures are also available, which conduct electricity and
heat up the glass surface to keep it fog free.
22.7.4 electronics
Carbon nanotubes have been used as conduits for electricity in very small electrical circuits due to their
superior electrical properties and absence of electromigration.
22.7.5 energy
The most advanced nanotechnology projects related to energy are: storage, conversion, manufacturing
improvements by reducing materials and process rates, energy saving (by better thermal insulation, for
example), and enhanced renewable energy sources. A reduction of energy consumption can be reached by
better insulation systems, by the use of more efficient lighting or combustion systems, and by use of lighter
and stronger materials in the transportation sector. Currently used light bulbs only convert approximately
5% of the electrical energy into light. Nanotechnological approaches like light emitting diodes (LEDs) or
quantum caged atoms (QCAs) could lead to a strong reduction of energy consumption for illumination.
22.7.6 space
Nanotechnology may hold the key to making space flight more practical. Advancements in nanomaterials
make lightweight solar sails and a cable for the space elevator possible. By significantly reducing the amount
of rocket fuel required, these advances could lower the cost of reaching orbit and travelling in space. Space
science, as long, played a role in the research and development of advancing technologies. Spacecraft are
being launched, with hulls that are composed of carbon fibres, a lightweight high strength material. Combine
that with smaller on–board computers that perform hundreds of times faster than computers used on spacecraft
just a decade ago, one can see the incredible advances in space exploration in just, past few years. The
advancements in material science and computer science have allowed the building, launching and deploying
of space exploration systems that continually do more and more as they become smaller and lighter.
22.7.6.1 Smart Materials
Some of the latest avenues being explored, which are more in the nano realm, in space science, include
smart materials for the hulls of spacecraft. These would be materials primarily composed of nanotube fibres
with nano sized computers integrated into them. These materials along with being even lighter will also be
far stronger too. One idea is to create a surface that will help transfer the aerodynamic forces working on a
spacecraft during launch. When the craft is launched the nano computers will flex the crafts hull to offset
pressure differences in the hull caused by the crafts acceleration through the atmosphere.
Then the same nano computer network in the hull would go to work heating the shaded side of the craft
and cooling the sun exposed side and to even create heat shielding for re-entry. To equalize the surface
temperature, now a spacecraft must be kept rotating and although a slight spin is good in maintaining the
attitude of a craft. Sometimes it interferes with the mission plan, like when a spacecraft is taking photographs
or is in the process of docking with another craft.
22.7.6.2 Swarms
Another avenue being investigated is a concept of nano robotics called “Swarms”. Swarms are nano robots
that act in unison like bees. They theoretically will act as a flexible cloth like material and are composed of
what’s called Bucky tubes. This cloth will be as strong as diamond. Add to this cloth of nano machines nano
computers and we have smart cloth. This smart cloth could be used to keep astronauts from bouncing around
inside their spacecraft while they sleep, a problem that arises when the autopilot computer fires the course
correction rockets. The cloth like material will be able to offset the sudden movements and slowly move the
sleeping astronaut back into position. Still another application for the nano robot swarms, being considered,
is that the smart cloth could be used in the astronauts’ space suits.
22.8 LimitatiONs Of NaNOtecHNOLOgy LO5

Any technology is generally known to solve a problem, but at the same time it may create a new problem or
has a limitation. There are potential risks of nanotechnology as well. A few limitations of the nanotechnology
are discussed below.
(i) The problem with nanotechnology is the scale it works on. Small and highly complex Nano machines
are probably awesome at building things with carbon only. Moreover, nanites may not work well in
the strong magnetic fields which are needed to produce permanent magnets. They may also not be
able to meet or operate under the conditions necessary to manufacture superconductors.
(ii) In view of constant breaking and reforming bonds as well as the friction caused by millions of
nanobots a large amount of heat is expected to be produced, which will lead problems in fabricating
certain extremely temperature sensitive materials. For example, nitroglycerin would explode if it
experienced even a one-degree shift in temperature. So large stations would likely be limited in their
production by the amount of heat they are able to disperse. Also, because the ratio of surface area
to internal space decreases as the size of an object increases, larger ships will have a harder time
dispersing heat than smaller vehicles.
(iii) It will be difficult for nanites to build something in a hostile environment. The workforce may be
billions of tiny robots strong but they will likely be unable to finish the job in a sandstorm or leave
large chunks of partially completed product scattered over a large area.
22.9 DisaDvaNtages Of NaNOtecHNOLOgy LO5

Nanotechnology has made significant impact in all fields of life and it has impressive applications in almost
every area of work, but it also has some major disadvantages which need to be addressed.
(i) Nanotech devices and machines have taken place of human to work faster and accurately. This has
lessened the importance of men power in the field of practical work, and reduced the employment
in the fields of traditional farming and manufacturing and industrial sector.
(ii) Nanotechnology has made atomic weapons more powerful and more destructive. However, this has
produced a big threat with regard to their easy accessibility. Unauthorized and criminal bodies can
reach nuclear weapons easily and its formulation could be stolen.
(iii) Diamond is now being produced massively with the help of nanotechnology, which has reduced the
value of diamond and increased the fall of diamond markets. Presence of alternative has decreased
the demand because alternates are more efficient and do not require the use of fossil fuels. The
manufacturer can now produce bulk of the products at molecular scale and decomposition is done to
create new components.
(iv) The presence of nanomaterials which contain nanoparticles is not in itself a threat but their increased
reactivity and mobility can make them risky. Nanotechnology has increased risk to the human
health, as the nanoparticles due to their small size can cause inhalation problem and many other fatal
diseases. Apart from what happens if non-degradable or slowly degradable nanoparticles accumulate
in organs, another concern is their potential interaction with biological processes inside the body.
(v) Nanotechnology is the most expensive technologies, whose cost is increasing day by day due to
the molecular structure and processing of the product. It has become difficult for the manufacturers
to randomly produce dynamic products due to the huge pricing of nanotech machines. This is
unaffordable for the common people.
(vi) Nanotechnology has raised the standard of living, but at the same time it has increased the pollution
including water and air pollutions due to the wastes generated by Nano devices or during the
nanomaterials manufacturing process. This pollution, called Nano pollution, may be very dangerous
for living organisms. Most of the human made nanoparticles do not appear in nature, so living
organism may not have appropriate means to deal with Nano waste. Hence, the whole life cycle of
these particles needs to be evaluated with respect to their fabrication, storage, distribution, potential
abuse and disposal. The impact on the environment may vary at different stages of the life cycle.
Concerns are raised about Nano pollution as it is not currently possible to precisely predict or
control ecological impacts of the release of these Nano products into the environment.
(vii) Finally, there are educational gap risk issues with regard to the nanotechnology, though it offers
rapid advances across many areas of science and engineering which are crucial to the society. For
example, the knowledge within scientific and industrial communities is not appropriately shared
with the civil society, public and regulatory agencies. Because of this innovative opportunities may
be lost and public confidence in transparency and accountability may wear away.
s UmmarY
Summary of the topics covered in this chapter is given below.

✦ Nanoscience is the study of objects having at least one dimension less than hundred nanometers; whereas
nanotechnology is the engineering of these objects using different techniques.
✦ Nanotechnology deals with various structures of matter having dimensions of the order of a billionth of
a meter. In any material, an atom on the surface is different from an atom in the bulk. This difference is
primarily because of the different potentials that the two atoms are experiencing.
✦ The ratio of number of atoms in the bulk to the number of atoms on the surface reveals that as the size of
the particle decreases the number of atoms on the surface is in a greater proportion of the total number
of atoms. Hence, as the particle size becomes smaller and smaller the surface atoms start dominating
the properties of the whole material.
✦ Nanoparticles are generally considered to be a number of atoms or molecules bonded together such that
the dimension of the bonded entity is of the order of 100 nm.
✦ The properties of nanomaterials are different from those of bulk because of two main reasons—the
surface effect and the quantum effect. The varied applications of nanoparticles is a consequence of their
varied properties.
✦ Synthesis of nanoparticles can be achieved by mechanical means such as ball milling; or by techniques
such as Gas condensation, Sputtering, Vaccum Deposition and Vaporization, Chemical Vapour Deposition
(CVD). Chemical Vapour Condensation (CVC), Sol-Gel Technique and Electrodeposition.
✦ Nanoparticles have properties different from those of bulk. Metallic nanoparticles exhibit Surface
Plasmon Resonance (SPR) which enables their use in medical diagnostics. In case of semiconductor
nanoparticles, the band gap increases with decrease in size leading to a change in absorption spectra
compared to bulk material.
✦ Quantum confinement is the restricted motion of randomly moving electrons in specific energy levels on
reduction of size. Quantum confinement leads to change in optical and electrical properties. Based on the
number of free dimensions available; nanostructures are classified as 3D, 2D, 1D or 0D. 3D structures
have none of its dimensions confined whereas 2D nanostructures have one dimension confined e.g. thin
films. 1D nanostructures have two dimensions confined e.g. nanowires and 0D nanostructures have all
their dimensions confined e.g. Quantum Dots.
✦ Nanowires are 1D structures and have an Aspect Ratio > 1000. They have diameter £ 10 nm and
unconstrained length. In Top down approach to Nanowire synthesis, etching and Lithographic techniques
are employed. For bottom-up approach, the Vapour Liquid Solid (VLS) method, Vapour Solid (VS)
Method and Catalyst free Template assisted methods are used.
✦ Carbon nanotubes, which are long and thin cylinders of carbon, were discovered in 1991 by S. Iijima.
These are large macromolecules that are unique for their size, shape, and remarkable physical properties.
They can be thought of as a sheet of graphite (a hexagonal lattice of carbon) rolled into a cylinder.
✦ Cartoon Nanotubes (CNTs) are single walled (SWCNT) or Multiwalled (MWCNT). MWCNTs have
higher strength than SWCNTs.
✦ Procedures used for CNT synthesis are: Arc Discharge method, Laser Ablation Method, Chemical
Vapour Deposition (CVD) method and n-hexane pyrolyris.
✦ CNTs have higher tensile strength than steels. They are highly elastic. The electrical conductivity of
CNTs is structure-dependent as the structure influences the collisions between conductive electrons and
atoms. The thermal conductivity of CNTs is 15 times superior to that of copper.
✦ Inorganic nanotubes are non-carbon nanotubes. They are cylindrical molecules that are often composed
of metal oxides or Group-III Nitrides.
✦ Biopolymers are polymers produced by living organisms. They are biodegradable. Biopolymers are
synthesized by chemical polymerization or by chemical modification of a naturally occurring polymer.
They find application in the packaging industry as food trays, wrappings, plastic etc.
✦ Fullerene is a carbon-molecule which could be in the form of a holow sphere, ellipsoid or tube. C-60
spherical fullerenes are known as Bucky balls. C-60, C-70, C-76, C-84, C-240 and C-540 are some of
the spherical fullerenes.
✦ The main methods of synthesizing fullerenes is by Electric Arc Discharge method, Laser Ablation
method and chemical Vapour Deposition method.
✦ Fullerenes are twice as hard as diamond when compressed to 70% of their original size. They are soluble
in organic solvents at room temperature.
✦ Quantum dots are 0-D nanostructures. The energy levels of Quantum dots are discrete, quantized and
isolated as in atoms. They find application in high power, low threshold semiconductor lasers, high
efficiency detectors and in solar cells.
✦ Commonly used techniques for fabrication of Quantum dots are Molecular Beam Epitaxy (MBE), Metal
Organic Chemical Vapour Deposition (MOCVD), and Pulsed Laser Deposition (PLD).
✦ Nanotechnology has got diverse applications in self-cleaning glass, clothing, scratch resistant coating,
electronics, energy, space and environment. Therefore, these applications of nanotechnology were
talked about.
s olVeD e XamPles
E XAMPLE 1 Estimate the Nb/Ns ratio for a spherical particle of diameter 12 mm and compare it with the same
for a nanoparticle of diameter 90 nm.
SOLUTION We have, Nb/Ns = r/3 where r is the radius of the particle. Hence, for the micrometer sized particle,
Nb/Ns = 6 ¥ 10–6/3 = 2 ¥ 10–6
Similarly for the nanometer sized particle,
Nb/Ns = 45 ¥ 10–9/3 = 15 ¥ 10–9
Hence, the ratio is smaller for the nanoparticle by a factor of 7.5 ¥ 10–3.
E XAMPLE 2 How is it possible to obtain nanoparticles of the same material yet having different colours?
SOLUTION In the case of metallic nanoparticles the color of the particle depends on its surface plasmon resonance (SPR)
frequency. Now it is known that the surface plasmon resonance frequency depends on the size as well as the shape of the
nanoparticle at hand and shifts as any of these parameters are changed. Hence, even for the same material it is possible to
have different SPR frequency for different particle sizes or shapes and thus a different colour.
E XAMPLE 3 How does a Single Electron Transistor (SET) work?

SOLUTION In a SET, the modulation is done by charge transfer on to a small island. This island has dimensions in the
nanometer regime and even the transfer of a single charge e to it can change its potential substantially. So we can put a
transistor in off state by transferring a single electron to the island and increasing its potential to a value so that the current
transfer between the source and drain ceases. Thus, it is given the name of a Single Electron Transistor.
E XAMPLE 4 What gives the increased yield strength in the systems that have a grain size in nano regime?
SOLUTION The reason for the increase in yield strength as the grain size becomes smaller and smaller is that smaller
grain sizes have more grain boundaries that offer resistance to the movement of dislocations. And since a material fails
when the dislocations gather at a spot and yield to breakage, by limiting movement of dislocations the yield strength is
increased.
Q.1 A special molecule of carbon made up of 60 carbon atoms is understood as a structure that shows
potential for a basic building block in the area of molecular manufacturing. The nontechnical name of
these molecule is
(a) Fullerene (b) Nano rods (c) Bucky balls (d) Nanotubes
Q.2 Graphene is
(a) A one-atom thick sheet of carbon
(b) A new material made from carbon nanotubes
(c) Thin film made from fullerene
(d) A software tool to measure and graphically represent nanoparticles
Q.3 Single-walled carbon nanotubes (SWCNTs) are
(a) Poor conductor (b) Excellent conductor
(c) Poor conductor than MWCNTs (d) None of the above
Q.4 Diameter of Bucky ball is about
(a) 100 Ao (b) 10 Ao (c) 1 Ao (d) 1000 Ao
Q.5 Properties of nanoparticles
(a) Are significantly different from the properties of bulk materials
(b) Are little bit different from the properties of bulk materials
(c) May be the same as in bulk material
(d) Are none of the above
Q.6 Surface area per unit volume for nanoparticles is
(a) Higher than macro-sized particles (b) Same as macro-sized particles
(c) Lower than macro-sized particles (d) None of the above
Q.7 Starch and cellulose are both biopolymers of
(a) Glucose (b) Maltose (c) Starch (d) Fructose
Q.8 Self-healing pain can be cured through
(a) Biotechnology (b) Information technology
(c) Nanotechnology (d) None of the above
Q.9 Carbon nanotubes are
(a) Hollow cylinders made up of carbon atoms
(b) Circular tubes made up of graphite
(c) Nanotubes made of carbon sheet
(d) Nothing but simple carbon atoms
Q.10 Upon decreasing the dimension of a nanoparticle what kind of a shift is observed in the absorption
spectra of a semiconducting particle?
(a) Red shift (b) Blue shift (c) Green shift (d) Violet shift
Q.11 What kind of a quantum mechanical system has a constant density of states?
(a) 1-D (b) 2-D (c) 3-D (d) 0-D
Q.12 A quantum mechanical system was found to have spikes in the plot of its density of states. Among the
following physical systems which can represent such a system?
(a) A quantum well (b) A bulk system
(c) A quantum dot (d) A carbon nanotube
Q.13 Which of these pairs correctly represents the constituent particles of an exciton?
(a) Electron, positron (b) Electron, hole
(c) Electron, positronium (d) Hole, hole
Q.14 Which among these systems can show us excitonic effects?
(a) Low T, low purity, bulk (b) Low T, low purity, quantum dot
(c) Room T, high purity, bulk (d) Room T, high purity, quantum dot
Q.15 Which of the following methods are not currently employed for fabricating carbon nanotubes?
(a) Arc discharge (b) Laser ablation
(c) Chemical vapour deposition (d) Ball milling
Q.16 Which of these does not represent a type of carbon nanotube?
(a) Chiral (b) Zigzag (c) Wavy (d) Armchair
Q.17 Which of the following is not a stage of Sol-gel formation?
(a) Agglomeration (b) Condensation (c) Hydrolysis (d) Sputtering
Q.18 Which of the following elements is known for its anti-bacterial properties in its nanoparticle form?
(a) Ag (b) Fe (c) Pd (d) Cu
Q.19 The absence of which of the following phenomenon is most critical in making carbon nanotubes as a
good conduit material in electronic circuits?
(a) Electromigration (b) Mechanical strain memory
(c) Poor mechanical strength (d) Thermal anisotropy
Q.20 Nanotechnology was brought into day light by delivering lectures by
(b) Max Planck (b) Einstein (c) Feynman (d) Lorentz
Q.21 Which of the following nano-particles are mostly used in paint industry
(a) Ag (b) Fe (b) TiO2 (c) SiO2
Q.22 Which of the following is correct for surface area to volume ratio in nanomaterials
(a) moderate (b) very less (c) very large (d) None of these
Q.23 The third known form of pure carbon is
(a) Diamond (b) Fullerene (c) Graphite (d) None of these
Q.24 When a bulk material is changed into nanoparticle which of the following will change state
(a) Physical (b) Chemical (c) both (a) & (b) (d) None of these
Q.25 Which of the following method is used to prepared carbon nanotube?
(a) Plasma arc-evaporation method (b) Chemical vapour deposition method
(c) Laser ablation method (d) All of these
Q.26 Bucky ball is the cluster of carbon atoms
(a) 10 (b) 60 (c) 15 (d) None of these
Q.27 Which of the following statement is correct for carbon nanotube
(a) almost 20 times stronger than steel (b) almost 6 times lighter than steel
(c) (a) & (b) (d) None of these
Q.28 Which of the nanocrystalline are used as separator plate in new generation batteries
(a) Nickel (b) Nickel hydrides (c) both (a) & (b) (d) None of these
Q.29 Carbon atom attain the shape in Bucky ball
(a) Hexagonal (b) pentagonal (c) trigonal (d) None of these
Q.30 Gold nanosphere of size 100 nm appears
(a) violet (b) red (c) orange (d) green
P ractice P roblems
General Questions
Q.1 Write a short note on nanotechnology.
Q.2 What do you understand by nanoparticles? Discuss their optical properties.
Q.3 Based on I-V characteristics, discuss the electrical properties of nanoparticles.
Q.4 In the light of mechanical properties of nanoparticles, explain how the yield strength varies with grain
size?
Q.5 What do you understand by quantum confinement? Discuss density of states for different types of
quantum confinements.
Q.6 What are carbon nanotubes? Discuss how various types of carbon nanotubes can be formed from
graphene?
Q.7 Discuss in short various techniques for the synthesis of nanoparticles.
Q.8 List the difference between chemical vapour deposition (CVD) and chemical vapour condensation (CVC).
Q.9 Explain various steps involved in sol-gel technique used for the synthesis of nanoparticles.
Q.10 Discuss how nanotechnology is useful in environment and space?
Q.11 How nanomaterials are different from bulk materials?
Q.12 Discuss the difference between nanoscience and nanotechnology.
Q.13 Discuss the basic difference between 0D, 1D, 2D and 3D materials.
Q.14 Write a note on nanowires and their different kinds.
Q.15 Give a brief description of synthesis techniques of nanowires.
Q.16 What do you understand by carbon nanotubes? These structures fall within which category?
Q.17 Discuss single-walled and multi-walled carbon nanotubes along with their differences.
Q.18 What are the methods for fabrication of carbon nanotubes? Discuss in brief.
Q.19 Discuss properties of CNTs.
Q.20 Write a note on inorganic nanotubes.
Q.21 Write a note on biopolymers.
Q.22 Discuss in brief 2D nanomaterials with respect to their synthesis.
Q.23 Discuss sol-gel technique used for the synthesis of nanoparticles.
Q.24 Write down the details of Ball milling and gas condenstion techniques used for the synthesis of nano-
particles.
Q.25 What are surface plasmons.
Q.26 What are Bucky balls or Fullerenes? How are these synthesized using electric arc discharge, laser
ablation and CVD techniques?
Q.27 Discuss in short the properties of fullerenes and their potential applications.
Q.28 What do you understand by quantum dots? How are these fabricated and characterized?
Q.29 Discuss any five applications of nanotechnology.
Q.30 What are the limitations of nanotechnology?
Q.31 Discuss five disadvantages of nanotechnology.
Appendices
Appendix
Measurements and errors 1

If you are asked to measure the same object two different times, there is always a possibility that the two
measurements may not be exactly the same. Then the difference between these two measurements is called
a variation or error in the measurements. In general, there are two types of errors in measurement. These are
called static error and dynamic error.
A1.1 StAtic Error And itS clASSificAtion

Static error of a measuring instrument refers to the numerical difference between the true value of a quality
and its measured value of quantity. This gives different indications when the same quantity is measured
repeatedly. On the other hand, dynamic error is the difference between the true value of a quantity changing
with time and the value indicated by the instrument.
Static errors are classified into three categories, namely, gross errors or human errors, systematic errors and
random errors.
A1.1.1 Gross Errors
Gross errors include all the human mistakes while reading and recording. Mistakes carried out in calculating
the errors also fall within this category. For example, while taking the reading from the meter of the instrument,
a person may read 21 as 27 or 31. Gross errors can be avoided if proper care is taken in reading, recording
the data and doing calculations accurately. We can also reduce such errors by increasing the number of
experimenters and by taking the average of more readings.
A1.1.2 Systematic Errors
Systematic errors are the errors which tend to be in one direction (either positive or negative). Systematic
errors include instrumental, environmental and personal errors. Instrumental error may be due to wrong
construction or calibration of the measuring instruments. These errors also include the loading effect, misuse of
the instruments and zero error in the instrument. Environmental error arises due to external conditions, which
include temperature, pressure, humidity, external magnetic field, etc. We can minimize the environmental
errors by maintaining the temperature and humidity of the laboratory constant through some arrangements,
and ensuring that there is no external magnetic or electrostatic field around the instrument. On the other hand,
personal errors are due to wrong observations, which may be due to lack of proper setting of the apparatus
or individual carelessness in taking observations.
A1.1.3 Random Errors
After calculating all the systematic errors, it is observed that there are still some more errors in the
measurement. These random errors are those errors which occur irregularly and are random with respect
to their sign and size. Random and unpredictable fluctuations in temperature, voltage supply and mechanical
vibrations of experimental set-up may lead to random errors. The important property of a random error is that
it adds variability to the data but does not affect average performance for the group. For this reason, random
error is sometimes referred to as noise.
A1.2 dynAmic Error

Dynamic errors are caused by the inertial properties of measuring instruments or equipment. Consider that
a varying quantity is recorded with the help of a recording device. Then the difference between the obtained
function and the actual process of change of the recorded quantity in time is called the dynamic error of
the given dynamic instrument. It is clear that these errors are caused by the time variation in the measured
quantity.
A1.3 SourcES of ErrorS

Sources of error, other than the inability of a piece of hardware to provide true measurements, are as under:
1. Insufficient knowledge of process parameters and design conditions
2. Errors caused by the person operating the instrument
3. Change in process parameters, irregularities, etc.
4. Poor maintenance
5. Poor design
6. Certain design limitations
A1.4 AccurAcy And PrEciSion

If you obtain a weight measurement of 6.5 kg for a given substance, and the actual or known weight is
10 kg, then your measurement is not accurate. It means your measurement is not close to the known value.
The closeness of a measured value to a standard or known value is referred to as accuracy. On the other hand,
precision refers to the closeness of two or more measurements to each other. The precision tells us to what
resolution (or limit) the quantity is measured. In the above example, if you weigh a given substance six times
and get 6.5 kg each time, then your measurement is very precise. It means the precision is independent of the
accuracy. You can be very precise but inaccurate, and also you can be accurate but imprecise. Moreover, you
can have accuracy without precision. For example, on an average your measurements for a given substance
may be close to the known value, but the measurements may be far from each other.
A1.5 rESolution
Resolution is the fineness to which an instrument can be read. We can take the example of two stopwatches,
out of which one is analog and the other is digital. Both are manually actuated and are looked at for resolution.
Appendix 1: Measurements and Errors 781
The analog stopwatch has to be viewed on its dial. If we look closely, we can relate the big hand to the
smallest tick mark on the big dial. That tick mark is a tenth of a second. It means the best a good eye can do is
resolve a reading to 1/10 second. Hence, this is the resolution of the stopwatch. On the other hand, the digital
stopwatch has two digits beyond the seconds. So it subdivides time in hundredths of a second. Since it is easy
to read to 1/100 of a second, the resolution of the digital stopwatch is 1/100 second.
A1.6 mEASurEmEnt uncErtAinty

Certainty is perfect knowledge which has total security from error.
Certainty is also the mental state of being without doubt. Every
experiment is approximate due to error in measurement. When a
number of measurements are done to a stable voltage (or other
parameter), the measured values will show a certain variation.
This variation is caused by thermal noise in the measuring circuit
of the measuring equipment and the measurement set-up. These
variations or the uncertainties are shown in Figure 1.1.
Figure 1.1
A1.7 StAndArd dEviAtion And vAriAncES

The uncertainty is estimated by calculating the standard deviation. Let x1, x2, x3 … xN be the results of an
experiment repeated N times. Then the standard deviation s is defined as
N
1
s=
N
Â ( xi - x )2
i =1
Ê 1 ˆ
where, x Á ∫
Ë N
Â xi ˜¯ represents the average of all the values of x. In this case, the uncertainty is of the
i
order of ±s. The standard deviation is defined in terms of the square of the deviations from the mean, which
N
is clear from the term Â ( xi - x )2 in the above formula. Moreover, s2 is known as the variance of the data.
i =1
The standard deviation s is the root mean square deviation of the data, measured from the mean.
A1.8 AbSolutE Error

Absolute error is defined as the magnitude of difference between the actual and the approximated values of
any quantity. For example, we measure a given quantity n times and a1, a2, a3, … an are the individual values.
Then the arithmetic mean (say am) can be found as
a1 + a2 +  + an
am = (i)
n
This can also be written as

n
Â ai
i =1
am = (ii)
n
The absolute error can now be obtained from the following formula:
Dan = am – an (iii)
In simple words, absolute error is the amount of physical error in a measurement. If we use a metre stick
to measure a given distance, then the error may be ±1 mm or ±0.001 m. This is the absolute error of the
measurement. For the measurement of a quantity x, the absolute error is Dx.
A1.9 rElAtivE Error

Relative error gives us an idea of how good a measurement is relative to the size of the object being measured.
For example, we measure the height of a room and the length of a small table by using a metre stick. We find
the height of the room as 3.256 m ± 1 mm and the length of the table as 0.082 m ± 1 mm. Then
Absolute error
Relative error = (i)
Value of thing measured
\ Relative error in measuring the height of the room
0.001 m
= = 0.0003
3.256 m
Relative error in measuring the length of the table
0.001 m
= = 0.0122
0.082 m
Here, it is clear that the relative error in measuring the length of the table is larger than the relative error in
measuring the height of the room. In both the cases, however, the absolute error is the same.
A1.10 PERcEntAGE ERRoR

If the relative error is represented in percent, then the error is called percentage error. For example, in the
above example, the percentage error in measuring the height of the room is 0.0003 ¥ 100 = 0.03%, while the
percentage error in measuring the length of the table is 0.0122 ¥ 100 = 1.22%.
The absolute error, relative error and percentage error can be summarized as follows. If we represent the
given or actual value by a and its approximated value as aapp, then
Absolute error = a - aapp
a - aapp aapp
Relative error = = 1-
a a
aapp
Percentage or percent error = 1 - ¥ 100
a
The important point is that the actual value a talked about cannot be zero.
A1.11 ERRoRS occuRRinG in ARithmEtic oPERAtionS

1.11.1 Addition and Subtraction
Let us consider two measured values a ± da and b ± db in which a and b are actual values, whereas da and db
are absolute errors of a and b, respectively.
The error obtained in the sum of these quantities is given by
Q ± dQ = (a ± da) + (b ± db)
= (a + b) + (±da ± db)
fi dQ = da + db (i)
Similarly, the error obtained in their difference is given by
Q ± dQ = (a ± da) – (b ± db)
= (a – b) + (±da ± db)
fi dQ = da + db (ii)
Hence, in arithmetic operations of addition and subtraction, the absolute error in the resultant is the sum of
the absolute errors of individual quantities. So errors are always added in these operations.
A1.11.2 multiplication and Division

Let us consider the two measured quantities a ± da and b ± db, where a and b are actual quantities and da and
db are the absolute errors in a and b, respectively. The errors occurring in multiplication and division can be
estimated as given below.
Q ± dQ = (a ± da) (b ± db) = ab ± bda ± adb ± dadb (iii)
On dividing by Q on the LHS and by ab on the RHS, we obtain
dQ da db
1± =1± ± (iv)
Q a b
In terms of percentage error, we have
dQ da db
¥ 100 = ¥ 100 + ¥ 100
Q a b
For general treatment, we can consider a quantity
Q = kax bycz (v)
Then error can be determined as
dQ da db dc
=x +y +z (vi)
Q a b c
It means, if Q is a measure of x power of a, y power of b and z power of c, then the fractional error will be the
sum of x times of fractional error in a, y times of fractional error in b and z times of fractional error in c. This
also infers that maximum error will be encountered due to the quantity carrying highest power. Therefore, in
experiments, that quantity should be determined with more precision in order to minimize the error.
S OLVED E XAMPLES
E xamplE 1 Find the standard deviation of the numbers 1, 5, 6, 7, 8, 10, 12.

Solution The mean of these numbers is found to be 7. The deviations are –6, –2, –1, 0, 1, 3, 5, respectively. Now using
the formula
N
1
s =
N
Â ( xi - x )2
i =1
1
= [36 + 4 + 1 + 0 + 1 + 9 + 25]
7
76
=
7
So, s = 10.86
= 3.295
E xamplE 2 If two resistances given as R1 = (50 ± 5) W and R2 = (150 ± 2) W are connected in series, then
find the equivalent resistance.
Solution R = (50 ± 5) + (150 ± 2)
= (50 + 150) ± (5 + 2)
= (200 ± 7) W
E xamplE 3 If the mass of a bulb with air is 98.625 ± 0.002 g and the mass of an empty bulb is 98.305 ± 0.002
g, then find the mass of air.
Solution Error in difference = (a ± da) – (b ± db)
= (a – b) + (±da ± db)
= (98.625 – 98.305) ± (0.002 + 0.002)
= 0.320 ± 0.004 g
E xamplE 4 If the capacity of a capacitor is C = 2 ± 0.4 F and the applied voltage is V = 20 ± 0.2 V, then find
the charge on the capacitor.
Solution Charge on capacitor, Q = CV = 2 ¥ 20 = 40 C
0.4
Percentage error in C = ¥ 100 = 20%
2
0.2
Percentage error in V = ¥ 100 = 1%
20
\ percentage of error in Q = 20 + 1 = 21%
21
or error in Q = 40 ¥ = 8.4 C
100
Hence, charge on the capacitor Q = 40 ± 8.4 C
E xamplE 5 The volumes of two bodies are measured to be V1 = (10.2 ± 0.02) cm3 and V2 = (6.4 ± 0.01) cm3.
Calculate the sum and difference in volumes with error limits.
Solution V1= (10.2 ± 0.02) cm3
V2 = (6.4 ± 0.01) cm3
DV = ± (DV1 + DV2)
= ± (0.02 + 0.01) cm3
= 0.03 cm3
E xamplE 6 The mass and density of a solid sphere are measured to be (12.4 ± 0.1) kg and (4.6 ± 0.2) kg/m3
respectively. Calculate the volume of the sphere with error limits.
Solution Here, m ± Dm = (12.4 ± 0.1) kg
r ± Dr = (4.6 ± 0.2) kg/m3
m 12.4
Volume V = = = 2.69 m3 = 2.7 m3
r 4.6
DV Ê Dm Dr ˆ
=±Á +
V Ë m r ˜¯
Ê Dm Dr ˆ
DV = ± Á + V
Ë m r ˜¯
Ê 0.1 0.2 ˆ
= ±Á + ¥ 2.7 = ± 0.14
Ë 12.4 4.6 ˜¯
V ± DV = (2.7 ± 0.14) m3
E xamplE 7 A current of 3.5 ± 0.5 A flows through a metallic conductor and a potential difference of
21 ± 1 volts is applied. Find the effective resistance of the wire.
Solution Given V = 21 ± 1 volts, DV = 1, I = 3.5 ± 0.5 A
DI = 0.5 A
V (21 ± 1)
Resistance R = = = 6.01 ± DR
I (3.5 ± 0.5)
DR DV DI
= error in measurement = +
R V I
1 0.5
= +
21 3.5
= 0.048 + 0.143 = 0.19
fi DR = 0.19 ¥ R = 0.19 ¥ 6 = 1.14
Effective resistance R = 6 ± 1.14 W
E xamplE 8 A rectangular board is measured with a scale having an accuracy of 0.2 cm. The length and
breadth are measured as 35.4 cm and 18.4 cm, respectively. Find the relative error and percentage error of
the area calculated.
Solution l = 35.4 cm, Dl = 0.2 cm
w = 18.4 cm and Dw = 0.2 cm
Area (A) = l ¥ w = 35.4 ¥ 18.4 = 651.36 cm2
DA Dl Dw
Relative error in area (dA) = = +
A l w
0.2 0.2
= + = 0.006 + 0.011 = 0.017
35.4 18.4
DA
Percentage error = ¥ 100 = 0.017 ¥ 100 = 1.7%
A
E xamplE 9 A physical quantity Q is related to four observables a, b, c, d as follows:
a3 b4
Q=
d2 c
The percentage errors of measurement in a, b, c and d are 1%, 3%, 4% and 3% respectively. What is the
percentage error in the quantity Q? If the value of Q calculated using the given relation is 8.768, to what value
should the result be rounded?
Solution Given
a3 b4
Q=
d2 c
Percentage error in Q is given by
DQ Da Db 1 Dc Dd
=3 +4 + +2
Q a b 2 c d
Da Db Dc Dd
Since = 1%, = 3%, = 4%, = 3%
a b c d
DQ 1
= 3 ¥ 1% + 2 ¥ 3% + ¥ 4% + 2 ¥ 3%
Q 2
= 3% + 6% + 2% + 6%
= 17%
\ percentage error in Q = 17%
If the calculated value of Q is 8.768, the roundoff value is 8.8.
E xamplE 10 Find absolute error, relative error and percentage error of the approximation 3.14 to the value p.
Solution Absolute error = 3.14 - p = 0.0015926536
3.14 - p
Relative error = = 0.000506957383
p
3.14 - p
Percentage error = ◊ 100% = 0.0506957383%
p
E xamplE 11 The refractive index (m) of water is found to have the values 1.29, 1.33, 1.34, 1.35, 1.32, 1.36,
1.30 and 1.33. Calculate the mean value, absolute error, the relative error and percentage error.
1.29 + 1.33 + 1.34 + 1.35 + 1.32 + 1.36 + 1.30 + 1.33
Solution mmean = = 1.3275  1.33
8
Absolute errors are
Dm1 = mmean – m1 = 1.33 – 1.29 = 0.04
Dm2 = mmean – m2 = 1.33 – 1.33 = 0.00
Dm3 = mmean – m3 = 1.33 – 1.34 = –0.01
Dm4 = mmean – m4 = 1.33 – 1.35 = –0.02
Dm5 = mmean – m5 = 1.33 – 1.32 = 0.01
Dm6 = mmean – m6 = 1.33 – 1.30 = 0.03
Dm7 = mmean – m7 = 1.33 – 1.33 = 0.00
Mean absolute error,

(| D m1 | + | D m2 | + | D m3 | + | D m 4 | + | D m5 | + | D m6 | + | D m7 | + | D m8 |
Dmmean =
8
= (0.04 + 0.00 + 0.01 + 0.02 + 0.01 + 0.03 + 0.03 + 0.00) ∏ 8
= 0.14 ∏ 8 = 0.0175  0.02
Relative error (dm) = Dmmean/mmean = 0.02 ∏ 1.33 = ±0.015  0.02
Percentage error = Dmmean/mmean ¥ 100% = ±0.015 ¥ 100 = ±1.5
l2 h
E xamplE 12 The radius of curvature of a concave mirror is given as R = + , where l and h are given as
6h 2
2 cm and 0.064 cm, respectively. Find the error in measuring the radius of curvature.
Solution l = 2 cm, Dl = 0.1 cm (LC of metre scale)
h = 0.064 cm, Dh = 0.001 cm (LC of spherometer)
l2 h DR 2 Dl -Dh Dh
R= + fi = + +
6h 2 R l h h
DR 2 Dl 2 Dh 2 ¥ 0.1 2 ¥ 0.001
= + = + = 00.1 + 0.03 = 0.131
R l h 2 0.064
E xamplE 13 The time of 30 oscillations of a simple pendulum whose length is 90 cm was observed to be 60
s. According to given data, find the value of g and determine percentage error in the value of g.
l l 60
Solution T = 2p fi g = 4p 2 ¥ 2 , T = = 2.00 s
g T 30
90 2
g = 4 ¥ 3.142 ¥ 2 = 887.364 cm/s
2
Maximum error in the value of g,
2l l 4p 2l
g = 4p 2 2
= 4p 2 2
= ¥ 302
T (t / 30) t2
Taking log on both sides,
loge g = loge 4 + 2 loge p + loge l – 2 loge t + 2 loge 30
loge g = 1.386 + 2.289 – 0.105 – 8.188 + 6.8
loge g = 2.182 or g = 8.86 m/s
Differentiating both sides, we get
Dg Dl Dt
= +2
g l t
l = 90 cm, Dl = 0.1 cm (LC of metre scale)

t = 60 sec, Dt = 0.1 sec (LC of stopwatch)
Dg 0.1 2 ¥ 0.1
= + = 0.0011 + 0.0033 = 0.0044
g 90 60
= 0.0044 ¥ 100% = 0.4%
E xamplE 14 In a measurement of the viscous drag force experienced by spherical particles in a liquid, the
force is found to be proportional to V1/3 where V is the measured volume of each particle. If V is measured to
be 30 mm3, with an uncertainty of 2.7 mm3, what will be the resulting relative percentage uncertainty in the
measured force? 2
Solution The relative percentage uncertainty in the measure of fore is 6 F 2 = ÊÁ ∂E ˆ˜ 6V 2
Ë ∂V ¯
Ê ∂E ˆ
6F = Á 6V , 6V Æ uncertainty in measurement of volume
Ë ∂V ˜¯
F µ V1/3
∂F 1 -2/3
µ V
∂V 3
1 1 1
fi 6F = ¥ 6V = 2 /3
¥ 2.7 = ¥ 2.7
3V 2/3
3(30) 3 ¥ 9.7
fi 6F = 0.09
E xamplE 15 One gram of salt is dissolved in water that is filled to a height of 5 cm in a beaker of 10 cm
diameter. The accuracy of length measurement is 0.01 cm while that of mass measurement is 0.01 mg. When
measuring the concentration c, what is the fractional error Dc/c?
Solution c = Mass/Volume
p d2
V = pr2h = h
4 2
2
Fractional error = Ê Dx ˆ + Ê Dy ˆ
ÁË x ˜¯ ÁË y ˜¯
2 2
DV Ê Dd ˆ Ê Dh ˆ Dd 0.01
= Á ˜ +Á ˜ = = 10 -3
V Ë d ¯ Ë h ¯ d 10
DV Dh 0.01
= 2 2 ¥ 10 -3 = - 2 ¥ 10 -3
V h 5
2 2
Dc Ê Dm ˆ Ê DV ˆ
= Á +Á = 10 -10 + (8 ¥ 10 -6 ) = .2 2 ¥ 10 -3 = 0.28%
c Ë m ˜¯ Ë V ˜¯
E xamplE 16 A battery powers two circuits C1 and C2 as shown in Figure 1.2.
Figure 1.2
The total current I drawn from the battery is estimated by measuring the currents I1 and I2 through the
individual circuits. If I1 and I2 are both 200 mA and the errors in the measurement are 3 mA and 4 mA
respectively, what is the error in the estimate of I?
Solution I1 = (200 ± 3) mA
I2 = (200 ± 4) mA
I = 400 ± DI
I = I1 + I2 = (400 ± 7) mA
DI = 7 mA
E xamplE 17 A resistance is measured by passing a current through it and measuring the resulting voltage
drop. If the voltmeter and ammeter have uncertainties of 3% and 4% respectively, then
(a) Find the uncertainty in resistance
(b) Find the uncertainty in the computed value of power dissipated across the resistance
Solution
(a) V = IR
Taking log on both sides and differentiating, we get
dV d I d R
= +
V I R
dR
±0.03 = ±0.04 +
R
dR
= ±0.07 (max.) = 7 %
R
(b) P = I2R
Taking log on both sides and differentiating, we have
d P 2d I d R
= +
P I R
= 2 ¥ 0.04 + 0.07
= 0.15 = 15%
Appendix
Optics 2
A2.1 ElEctromagnEtic (Em) SpEctrum
A large number of frequencies of electromagnetic waves visualized in numerical order constitutes an
electromagnetic (EM) spectrum. Frequencies that are usable for radio communication occur near the lower
end of the EM spectrum. With the increase of frequencies, the EM energy becomes dangerous to human
beings. For example, a microwave oven can be a hazard if it is not shielded properly. Also, with the increase
of frequencies, it becomes difficult to employ EM energy for communication.
The electromagnetic spectrum as per frequency range and usage of EM energy is given below in Table 2.1.
Table 2.1
Approximate Frequency Range EM Phenomena Examples of Uses
530–1600 kHz Radio waves AM radio
3–30 MHz Shortwave radio
50–250 MHz FM radio, VHF TV
450–800 MHz UHF TV
3–300 GHz Microwaves Radar, Satellite communication
3 4
10 –10 GHz Infrared radiation Photography
105–106 GHz Visible light Human vision
6 8
10 –10 GHz Ultraviolet radiation Sterilization
8 9
10 –10 GHz X-rays X-ray (medical)
10 13
10 –10 GHz g-rays Cancer therapy
14
> 10 GHz Cosmic rays Astronomy (Physics)
A2.2 lEnS abErrationS

The inability of a lens to form the white point image of an object is called lens aberration.
a2.2.1 Chromatic Aberration and its Kinds

When the rays of white light parallel to the principal axis are incident on a lens, the different colours of the
white light are refracted by different amounts at different places. It can also be said that the rays are dispersed
Appendix 2: Optics 791
into various colours and are focused at different distances from the lens (see Figure 2.1). This happens due to
the fact that the refractive index of the lens depends on the wavelength of light.
Actually, violet light gets refracted more than red light. Hence, the point at which violet light focuses is nearer
the lens than the point at which the red light focuses. Thus, the image formed by a lens is usually coloured
and blurred. This inability of a lens to form a single image of a white object is called chromatic aberration.
Figure 2.1
Depending upon the blurredness of the image along the axis or transverse (lateral) to the axis, the chromatic
aberration is of two types: longitudinal or axial chromatic aberration and lateral chromatic aberration.
Longitudinal Chromatic Aberration
The longitudinal aberration or axial chromatic aberration is actually the spreading of an image along the
principal axis. It means the longitudinal aberration is the formation of images of different colours at different
positions along the axis. The axial distance between the positions of red and violet images is a measure of
axial aberration. In Figure 2.1, (XR – XV) is the measure of this aberration. Clearly, the longitudinal aberration
is positive for the case of a convex lens and is negative for a concave lens.
Lateral Chromatic Aberration

If f is the focal length of the lens and u is the distance of an object from the lens, then the magnification is
given by
f
m=
u+ f
Since the magnification depends on the focal length of the lens, which is different for different colours, the
images of different colours will be of different sizes. This happens when the finite-sized white object is
placed on the axis of the lens. Based on this observation, we can define lateral chromatic aberration as the
formation of images of different sizes for different wavelengths due to variation of lateral magnification with
the wavelength.
A2.3 corrEction of chromatic abErrationS: uSE of DoublEt

For minimizing or reducing chromatic aberration, we make use of two lenses (doublet). This can happen
if two or more lenses are combined together in such a way that this lens combination produces images of
different colours at the same place or position and also of the same size. For this, a combination can be
prepared by using a crown-glass convex lens of low focal length and a flint-glass concave lens of greater focal
length. Such a combination of lenses is called an achromatic doublet, and the reduction or minimization of
chromatic aberration is termed as chromatism of two lenses in contact.
a2.3.1 theory of achromatism
Consider two lenses of different materials which are placed in contact with each other. Under this situation,
if the focal length can be found to be independent of the colours under certain conditions, the combination of
the two lenses will be called achromatic. The focal length f of a thin lens is given by
1 Ê 1 1 ˆ
= ( m - 1) Á - (i)
f Ë R1 R2 ˜¯
Ê 1 1 ˆ 1
or ÁË R - R ˜¯ = f ( m - 1) (ii)
1 2
Here, m is refractive index of the lens and R1 and R2 are the radii of curvature of the two surfaces of the lens.
If we take d f as the change in the focal length f, corresponding to a change dm in the refractive index m, then
df Ê 1 1 ˆ
one can find - 2 = dm Á - after differentiating Eq. (i).
f Ë R1 R2 ˜¯
By making use of Eq. (ii), this can be written as
df dm 1
- = ◊ (iii)
f2 ( m - 1) f
dm
The ratio (= w ) is the dispersive power of the lens between the two colours for which the difference
( m - 1)
in refractive index is dm and the mean refractive index is m. Hence,
Appendix 2: Optics 793
df w
- 2
= (iv)
f f
If f1 and f2 are taken as the mean focal lengths of two thin lenses of combination and w1 and w2 are the
dispersive powers between two colours for which the combination is to be achromatized, then the focal length
of the combination can be written as
1 1 1
= + (v)
F f1 f2
From this, we get
dF d f1 d f2
- 2
=- -
F f12 f22
d f1 w1 df w
In view of Eq. (iv), we can write - = and - 22 = 2
f12 f1 f2 f2
dF w1 w 2
\ - 2
= +
F f1 f2
The lens combination will be said to be achromatic if F does not change with colour, i.e., dF = 0. This yields
w1 w 2 f1 w
+ = 0 or =- 1 (vi)
f1 f2 f2 w2
So this condition is the required condition for a doublet. The negative sign
indicates that the combination should be a convex lens and a concave lens
(Figure 2.2). Crown Flint
This condition also conveys that the ratio of the focal lengths of the two
lenses is numerically equal to the ratio of dispersive powers of their materials.
Since w1 and w2 are positive quantities, the focal lengths must carry opposite Figure 2.2
signs, which justifies the combination of a convex lens and a concave lens
(Figure 2.2).
Appendix
Mechanical Properties of
Materials
3
A3.1 ElAsticity
Elasticity is a fundamental property of materials, and any material or body can be deformed by the application
of an external force. If the body returns to its original shape after the removal of the force, it is said to be
elastic. Springs of all kinds are examples of elastic bodies. Most substances are found to be elastic to some
degree. In technical terms, a substance with high elasticity is the one that requires a large force to produce a
distortion. For example, a steel sphere is a substance of high elasticity.
A3.2 strEss And strAin
We define certain terms such as stress and strain, for comparing the elasticity of materials. Consider a steel
wire, which is held rigidly at the top end and has a load fastened to the lower end. The wire under this
situation is said to be under stress, the magnitude of which is equal to the ratio of the applied force (the weight
in this case) to the cross-sectional area, i.e.,
F
Stress =
A
From this, one can observe that the SI unit of stress is N/m2.
If the load is significantly enhanced, the wire will be stretched by an amount DL (for its length as L). Under
this situation, we define another term called strain, which is a measure of the distortion of an object. Strain
is defined as the change in a spatial variable divided by the original value of that variable. If we take this
variable as the length, then
DL
Strain =
L
From this, one can observe that strain is a dimensionless quantity. It means it does not carry units.
There are three ways in which a body may change its dimensions under the action of an external force.
Consider a solid cylinder, which is stretched by two equal forces applied normal to its cross-sectional area.
The restoring force per unit area in this case is called tensile stress. Conversely, if the cylinder is compressed
under the action of applied forces, the restoring force per unit area is called compressive stress. Since in both
the cases, there is a change in the length of the cylinder, the tensile or compressive stress can also be termed as
longitudinal stress. The change in the length DL to the original length L of the cylinder, in this case, is known
Appendix 3: Mechanical Properties of Materials 795
as longitudinal strain. On the other hand. if two equal and opposite deforming forces are applied parallel to
the cross-sectional area of the cylinder, there would be a relative displacement (say Dx) between the opposite
faces of the cylinder. In this situation, the restoring force per unit area developed due to the applied tangential
force is known as tangential stress or shearing stress. The strain so produced by the tangential force is known
as shearing strain. This is defined as the ratio of relative displacement of the faces Dx to the length of the
cylinder L, whereas the volume strain is defined as the ratio of change in volume (DV) to the original volume
(V).
A3.3 HookE’s lAw

Hooke’s law is related to the stress and strain produced in a body under the action of an external force. This
law is an empirical law which is found to be valid for most materials. Hooke’s law states that the stress and
strain are proportional to each other for small deformations in the body. It means
Stress = k¢ ¥ Strain
where k¢ is the proportionality constant and is known as the modulus of elasticity.
A3.4 young’s Modulus
Young’s modulus is denoted by the symbol Y, and is defined as the ratio of tensile (or compressive) stress (s)
to the longitudinal strain (e). It means
s
Y=
e
Putting the values of s and e in the above formula, we get
( F /A) (F ¥ L)
Y= =
( DL /L ) ( A ¥ DL )
A3.5 sHEAr Modulus

The shear modulus of a material is denoted by G, and is defined as the ratio of shearing stress to the
corresponding shearing strain. The shear modulus is also called the modulus of rigidity, which is given by
( F /A)
G=
( Dx /L )
For small values of Dx, the ratio Dx/L can be defined as the angle q (shown in Figure 3.1). Then
F
G=
(A ¥ q)
Since F/A represents the shearing stress ss, we can have
ss = G ¥ q
It can be seen that the SI unit of shear modulus is N/m2.
Figure 3.1
A3.6 Bulk Modulus

Bulk modulus is denoted by B, and is defined as the ratio of stress to the volumetric strain. In view of the unit
of stress as N/m2, the bulk modulus can be defined in terms of pressure, p, as
p
B= -
( DV /V )
Since there is a decrease in volume with the increase of pressure, if p is positive and DV is negative, the
negative sign appears in the above formula. Thus, for a system in equilibrium, the value of the bulk modulus
B is always positive. The SI unit of bulk modulus is the same as that of pressure, i.e., N/m2 or Pa.
The reciprocal of the bulk modulus is known as compressibility, which is denoted by k. The compressibility
is basically the fractional change in volume per unit increase in pressure, i.e.,
Ê 1ˆ Ê 1 ˆ Ê DV ˆ
k = Á ˜ = -Á ˜ ¥ Á
Ë B¯ Ë Dp ¯ Ë V ˜¯
A3.7 rElAtion BEtwEEn ElAstic constAnts

Young’s modulus (Y), bulk modulus (B) and modulus of rigidity (G) are the elastic constants. These are
related to each other through Poisson’s ratio (s) as follows:
Y = 3B(1 – 2s) = 2G(1 + s)
These are also related to each other as per the relation
9 BG
Y=
(G + 3 B)
Relation 1: Y = 3B(1 – 2s)

We consider a cube A1B1C1D1A2B2C2D2 (Figure 3.2).
On each of its faces, an equal amount of force F is applied. Due to this action, the edges A1B1, B1A2 and B1C1
are elongated (in the direction of force), and are also contracted (in the direction perpendicular to the force).
If a represents the increase in length per unit length

per unit tension along the direction of the force, and b
represents the contraction produced per unit length per
unit tension in a direction perpendicular to the force,
then the Poisson’s ratio s = b/a. In view of this, the
elongation produced in the edges A1B1, B1A2 and B1C1
will be Fa each. Similarly, the contraction produced in
these edges will be 2Fb each, as the two pairs of forces
are acting perpendicular to each edge.
In view of the above, the lengths of the edges would be
A1B1 = 1 + Fa – 2Fb = B1A2 = B1C1.
Hence, the new volume of the cube will be
= (1 + Fa – 2Fb)3
= (1 + Fa – 2Fb) (1 + F2a2 + 2Fa
+ 4F2b2 – 4 Fb – 4 F2ab) Figure 3.2
Since a and b are small quantities, the terms F2a2, F2b2 and F2ab can be neglected. Thus, the new volume
becomes
(1 + Fa – 2Fb) (1 + 2Fa – 4Fb) = 1 + 2Fa – 4Fb + Fa – 2Fb
= 1 + 3F(a – 2b)
From this, we find the change in volume DV as
DV = 1 + 3F (a – 2b) – 1 = 3F(a – 2b) (i)
If a pressure P is applied on all the faces to compress the cube, it can be found that
DV = 3P(a – 2b) (ii)
From the above, the volumetric strain can be obtained as
DV/V = 3P (a – 2b) ∵V=1
Now, bulk modulus
Stress P
B= =
Volumetric strain 3P(a - 2 b )
1
or B=
3(a - 2 b )
Since Young’s modulus Y = 1/a and Poisson’s ratio s = b/a, we divide RHS of the above relation by a
fi 1/a Y
B= =
3(1 - 2b /a ) 3(1 - 2s )
This gives Y = 3 B(1 - 2s )

Relation 2: Y = 2G(1 + s)
We consider that a tangential force F is applied to the face
A1B1A2B2 of the cube as shown in Figure 3.3. Due to the action
of the force F, the face A1B1C1D1 gets displaced to A1¢B1¢C1D1
and the diagonal D1B1 gets elongated to D1B1¢ while the
diagonal C1A1 is decreased to C1A1¢.
In this case, the shearing stress along A1B1 will have the
same impact as the tensile stress along D1B1 and an equal
compression stress along C1A1 at right angles. If longitudinal
and lateral strains per unit stress are represented by a and b
respectively, then extension along D1B1 due to tensile stress
= D1B1 ¥ Shearing stress ¥ a (i)
Figure 3.3
Extension along D1B1 due to compression stress along A1C1
= D1B1 ¥ Shearing stress ¥ b (ii)
By adding Eqs (i) and (ii), we get the total extension along D1B1 as
= D1B1 ¥ Shearing stress ¥ (a + b) (iii)
Now, if q is very small, then –A1B1¢C1 ª 90° and –B1B1¢ N = 45°
Hence, the increase in length of D1B1 = B¢1N
= B1B¢1 ◊ cos 45°
B1 B¢1
= (iv)
2
From Eqs (iii) and (iv), we have
B1 B¢1
D1B1 ¥ Shearing stress ¥ (a + b) =
2
A1 B1 1
or Shearing stress ¥ = (∵ D1 B1 = B1 A1 2 )
B1 B1¢ 2(a + b )
Shearing stress 1
or =
(B1 B1¢ / A1 B1 ) 2(a + b )
Shearing stress 1
or =
Shearing strain 2(a + b )
1 1/ a
or G= =
2(a + b ) 2(1 + b / a )
Y
or G=
2(1 + s )
or Y = 2G (1 + s )
Relation 3: Y = 9 BG/(G + 3B)

It is clear from this relation that the Poisson’s ratio s does not appear here. Hence, we intend to eliminate s
from the following two relations:
Y = 3B(1 – 2s) (i)
Y = 2G (1 + s) (ii)
Equation (i) gives
1 – 2s = Y/3B (iii)
Equation (ii) gives
1 + s = Y/2G
or 2 + 2s = Y/G (iv)
Adding Eqs (iii) and (iv), we obtain
Y Y
+ =3
3B G
(G + 3 B)
or Y =3
3 BG
9 BG
or Y=
(G + 3 B)
A3.8 BEAMs And tHEir BEnding

A beam is a structural member that is acted upon by a system of external loads at right angles to the axis. These
loads acing perpendicular to the axis of the beam can deform the beam, which is called bending of beam. If
the plane of loading passes through one of the principal centriodal axes of the cross section of the beam, the
bending is said to be plane; otherwise the bending is said to be oblique. If the load can be considered to act
at a point, then the load is called a point load. On the other hand, if the load is distributed or spread in some
manner over the length of the beam, then it is called a distributed load. If the spread of the load is uniform, it
is said to be uniformly distributed; otherwise it is said to be a non-uniformly distributed load.
A3.8.1 Classification of Beams
Depending upon the type of support, beams are classified as cantilever, simply supported beam, overhanging
beam, fixed beam or continuous beam. A cantilever is a beam whose one end (say A) is fixed and the other
end (say B) is free. The length between A and B is known as the length of the cantilever. If both the ends of a
beam freely rest on a wall, columns or knife edges, the beam is called a simply supported beam. In all such
cases, the reactions are always upwards. If the supports are not situated at the ends of the beam, i.e., one or
both the ends project beyond the supports, then the beam is called an overhanging beam. On the other hand,
if both the ends of the beam are rigidly fixed or built-in into its supporting walls, then the beam is known as
a fixed beam. Finally, a continuous beam is defined as the beam that has more than two supports. In this case,
the supports at the extremes are called end supports and all the other supports, except the extreme, are called
intermediate supports.
A3.9 SheAring ForCe (SF) And Bending MoMent (BM)

With respect to the bending of beams, two terms,
i.e., shearing force (SF) and bending moment (BM)
are important to discuss. Consider a beam, which is
in equilibrium under the action of a series of forces.
When the beam is cut in some section XX (see
Figure 3.4) and its left part to the section remains in
equilibrium, then it is obvious that some force must
act at the section. This force would be provided by the
adjacent material prior to the cutting of the beam and
would act tangentially to the section. Hence, this is
clear that there will be shearing force at the section.
In view of the equilibrium of the material to the left of
the section XX, there can be no resultant moment to the Figure 3.4
left of the section. It means any moment produced by
the forces acting on the beam must be balanced by an equal and opposite moment produced by the internal
forces acting in the beam at the section. This is called the bending moment at the section, which is the
algebraic sum of moments to the left or right of the section. With respect to the equilibrium either for forces
or moments, the resultant caused by the applied forces to one side of the section is balanced by the bending
moment and shearing force acting at the section.
A3.9.1 Sign Convention
In general, a shearing force having an upward direction to the RHS of a section or downwards to the LHS of
the section is taken to be +ve. The shearing force having a downward direction to the right of the section or
upward direction to the left of the section will be –ve. Following the same terminology, a bending moment
causing concavity upwards is taken to be +ve and called sagging bending moment. On the other hand, a
bending moment causing convexity upwards is taken as –ve and called hogging bending moment.
A3.10 cAntilEvEr witH loAd At FrEE End

As mentioned earlier, a cantilever is supported at one end only and being built-in at its support, it gives a fixed
slope at that point. Now, we consider a section XX at a distance x from the free end B [see Figure 3.5(a)]. Then
shearing force at X = Sx = –M, and
bending moment at X = Bx = –Mx.
This shows that the shearing force is constant at all sections of the member between A and B. However,
the bending moment at any section is proportional to the distance of the section from the free end. The SF
diagram is shown in Figure 3.5(b) whereas the BM diagram is depicted in Figure 3.5(c).
X M
A x
B
X
L
(a)
M –ve
(b)
–ve
ML
(c)
Figure 3.5
A3.11 cAntilEvEr witH uniForMly distriButEd loAd

As is clear, here we consider that the load is
uniformly distributed over the whole length
of the beam. So it is appropriate to take the
loading as W per unit run [shown in Figure
3.6(a)].
Now, we consider the section XX at a distance
x from the free end B. In this case, the shearing
force at X = Sx = –Wx whereas the bending
moment at
Wx 2
X = Bx = Ú Sx dx = -
2
This shows that the variation of the shear
force is according to the linear law. However,
the variation of bending moment is parabolic,
being proportional to x2.
The shearing-force diagram is shown in Figure
Figure 3.6
3.6(b) and the bending-moment diagram is
given in Figure 3.6(c).
A3.12 cAntilEvEr witH uniForMly distriButEd loAd And A concEntrAtEd

loAd At tHE FrEE End
We consider a cantilever AB of length L, which carries a uniformly distributed load W per unit run over the
whole length and a concentrated load M at the free end. We take a section XX at a distance x from the free
end B [Figure 3.7(a)].
In this case, the shearing force and bending moment at the section X are given by
Ê Wx 2 ˆ
Sx = - ( M + Wx ) and Bx = Ú S x dx = - Á Mx +
Ë 2 ˜¯
Here also, we find that the shearing force follows the linear law, whereas the bending moment follows the
parabolic law. The corresponding variations are given in Figure 3.7(b) and Figure 3.7(c), respectively, for the
shearing force and bending moment.
Figure 3.7
S OLVED E XAMPLES
E xamplE 1 A structural steel rod has a radius of 10 mm and a length of 1.0 m. A 100 kN force stretches
it along its length. Calculate (a) stress, (b) elongation, and (c) strain on the rod. Young’s modulus (Y) of
structural steel is 2.0 ¥ 1011 Nm–2.
Solution We assume that the rod is held by a clamp at one end and the force F is applied at the other end, parallel to the
length of the rod. Then the stress on the rod is given by
F F 100 ¥ 103 N
Stress = = 2 = = 3.18 ¥ 108 Nm -2
A pr 3.14 ¥ 10 -4 m 2
The elongation
( F /A) L 3.18 ¥ 108 Nm -2 ¥ 1 m
DL = =
Y 2 ¥ 1011 Nm -2
= 1.59 ¥ 10–3 m = 1.59 mm
The strain is given by
Strain = DL/L = (1.59 ¥ 10–3 m)/1 m
= 1.59 ¥ 10–3 = 0.16 %
E xamplE 2 A square load slab of 50 cm side and 10 cm thickness is subject to a shearing force (on its narrow
force) of 9.0 ¥ 104 N. The lower edge is reverted to the floor (Figure 3.8). By how much will the upper edge
be displaced?
Solution The load slab is fixed and the force is applied parallel to the narrow force. The area of the force parallel to
which this force is applied is
A = 50 cm ¥ 10 cm = 0.5 m ¥ 0.1 m
= 0.05 m2
Therefore, the stress applied is
= (9.0 ¥ 104 N/0.05 m2) = 1.80 ¥ 106 Nm–2
50 cm
Figure 3.8
We know that
Dx Stress
Shearing strain = =
L G
Therefore, the displacement

Dx = (Stress × L)/G
= (1.8 ¥ 106 Nm–2 ¥ 0.5 m)/(5.6 ¥ 109 Nm–2)
= 1.6 ¥ 10–4 m = 0.16 mm
E xamplE 3 The average depth of the Indian Ocean is about 3000 m. Calculate the fractional compression DV/V of water
at the bottom of the ocean, given that the bulk modulus of water is 2.2 ¥ 109 Nm–2. (Take g = 10 ms–2).
Solution The pressure exerted by a 3000 m column of water on the bottom layer
p = hrg = 3000 m ¥ 1000 kg m–3 ¥ 10 ms–2
= 3 ¥ 107 kg m–1 s–2
= 3 ¥ 107 N m–2
Fractional compression,
DV Stress
=
V B
3 ¥ 107
=
2.2 ¥ 109
= 1.36 ¥ 10–2
or = 1.36 %
Chapter-wise Answers to
Objective Type Questions
Chapter 1
Q1. (b) Q2. (c) Q3. (c) Q4. (c) Q5. (c) Q6. (a) Q7. (d) Q8. (c)
Q9. (b) Q10. (c) Q11. (c) Q12. (c) Q13. (d) Q14. (d) Q15. (b) Q16. (b)
Q17. (c) Q18. (a)
Chapter 2
Q1. (a) Q2. (c) Q3. (b) Q4. (b) Q5. (a) Q6. (d) Q7. (b) Q8. (b)
Q9. (d) Q10. (d) Q11. (b) Q12. (a, c) Q13. (d) Q14. (a) Q15. (d) Q16. (a)
Q17. (a) Q18. (d)
Chapter 3
Q1. (c) Q2. (a) Q3. (d) Q4. (a) Q5. (c) Q6. (a) Q7. (a) Q8. (d)
Q9. (c) Q10. (a) Q11. (a) Q12. (b) Q13. (d) Q14. (d) Q15. (a) Q16. (a)
Q17. (a) Q18. (a) Q19. (c)
Chapter 4
Q1. (b) Q2. (a) Q3. (d) Q4. (a) Q5. (c) Q6. (a) Q7. (a) Q8. (b)
Q9. (c) Q10. (b) Q11. (b) Q12. (b) Q13. (a) Q14. (c) Q15. (a) Q16. (d)
Q17. (a) Q18. (a) Q19. (c) Q20. (d) Q21. (a)
Chapter 5
Q1. (b) Q2. (b) Q3. (d) Q4. (a) Q5. (a) Q6. (a) Q7. (b ) Q8. (d)
Q9. (d) Q10. (d) Q11. (a) Q12. (a)
Chapter 6
Q1. (c) Q2. (c) Q3. (a) Q4. (d) Q5. (b) Q6. (a) Q7. (a) Q8. (a)
Q9. (c) Q10. (b)
Chapter 7
Q1. (a) Q2. (d) Q3. (d) Q4. (d) Q5. (b) Q6. (b) Q7. (d) Q8. (b)
Q9. (b) Q10. (b)
Chapter 8
Q1. (a) Q2. (c) Q3. (b) Q4. (d) Q5. (c,d) Q6. (a,c) Q7. (b) Q8. (c)
Q9. (b) Q10. (c) Q11. (c) Q12. (a) Q13. (c) Q14. (b) Q15. (d) Q16. (b)
Q17. (a) Q18. (a) Q19. (b) Q20. (d)
Chapter 9
Q1. (d) Q2. (d) Q3. (b) Q4. (a) Q5. (a) Q6. (c) Q7. (b) Q8. (c)
Q9. (a) Q10. (b) Q11. (b) Q12. (a) Q13. (d) Q14. (a) Q15. (c) Q16. (c)
Q17. (b) Q18. (a) Q19. (a) Q20. (b) Q21. (b) Q22. (b) Q23. (c) Q24. (a)
Q25. (a) Q26. (d) Q27. (d) Q28. (d) Q29. (d) Q30. (b) Q31. (a) Q32. (d)
Q33. (d) Q34. (b) Q35. (c) Q36. (d) Q37. (d) Q38. (a) Q39. (b) Q40. (a)
Chapter 10
Q1. (b) Q2. (a) Q3. (c) Q4. (d) Q5. (a) Q6. (c) Q7. (a) Q8. (d)
Q9. (a) Q10. (a) Q11. (d) Q12. (a) Q13. (d) Q14. (b) Q15. (b) Q16. (b)
Q17. (a)
Chapter 11
Q1. (b) if the surface f is constant
Q2. (c) Q3. (c) Q4. (c) Q5. (a) Q6. (b) Q7. (a) Q8. (d) Q9. (c)
Q10. (a) Q11. (c) Q12. (d) Q13. (a) Q14. (a) Q15. (c) Q16. (c) Q17. (a)
Q18. (c) Q19. (b) Q20. (a) Q21. (b) Q22. (c) Q23. (d) Q24. (b) Q25. (d)
Q26. (c) Q27. (c) Q28. (b)
Chapter 12
Q1. (a) Q2. (b) Q3. (a) Q4. (b) Q5. (d) Q6. (b) Q7. (d) Q8. (a)
Q9. (d) Q10. (a) Q11. (b) Q12. (b) Q13. (b) Q14. (b)
Chapter 13
Q1. (b) Q2. (b) Q3. (c) Q4. (a) Q5. (a) Q6. (b) Q7. (b) Q8. (a)
Q9. (b) Q10. (a) Q11. (d) Q12. (a) Q13. (d) Q14. (d) Q15. (c) Q16. (a)
Q17. (d) Q18. (b) Q19. (c) Q20. (c) Q21. (c) Q22. (b) Q23. (b) Q24. (d)
Q25. (d) Q26. (a) Q27. (c) Q28. (d) Q29. (c) Q30. (a) Q31. (c) Q32. (a)
Q33. (a) Q34. (a) Q35. (a) Q36. (b) Q37. (a) Q38. (a) Q39. (a) Q40. (b)
Q41. (d) Q42. (a) Q43. (b) Q44. (a) Q45. (a) Q46. (a) Q47. (a) Q48. (c)
Q49. (b) Q50. (c) Q51. (a) Q52. (c) Q53. (a) Q54. (c) Q55. (a)
Chapter 14
Q1. (d) Q2. (c) Q3. (b) Q4. (d) Q5. (d) Q6. (c) Q7. (d) Q8. (b)
Q9. (a) Q10. (c) Q11. (d) Q12. (b) Q13. (b) Q14. (b) Q15. (a) Q16. (b)
Q17. (c) Q18. (b) Q19. (d) Q20. (a) Q21. (a) Q22. (a) Q23. (a) Q24. (b)
Q25. (c) Q26. (a) Q27. (b) Q28 (b) Q29. (a)
Chapter 15
Q1. (d) Q2. (d) Q3. (a) Q4. (a) Q5. (d) Q6. (d) Q7. (a) Q8. (d)
Q9. (c) Q10. (b) Q11. (b) Q12. (a) Q13. (d) Q14. (a) Q15. (b) Q16. (b)
Q17. (a) Q18. (a) Q19. (b) Q20. (b)
Chapter 16
Q1. (d) Q2. (b) Q3. (a) Q4. (b) Q5. (a) Q6. (d) Q7. (a) Q8. (a)
Q9. (b) Q10. (a) Q11. (d) Q12. (d) Q13. (b) Q14. (a) Q15 .(c Q16. (a)
Q17. (c)
Chapter 17
Q1. (b) Q2. (d) Q3. (c) Q4. (a) Q5. (d) Q6. (c) Q7. (a) Q8. (a)
Q9. (a) Q10. (c) Q11. (a) Q12. (d) Q13. (c) Q14. (b) Q15. (c)
Chapter 18
Q1. (c) Q2. (d) Q3. (b) Q4. (c) Q5. (d) Q6. (b) Q7. (b) Q8. (b)
Q9. (d) Q10. (b) Q11 (a) Q12. (c) Q13. (d) Q14. (c) Q9. (b) Q15. (b)
Q16. (b) Q17. (c) Q11. (c) Q18. (b) Q19. (a) Q20. (a) Q21. () Q22. (a)
Q23. (b) Q24. (a) Q25. (b) Q26. (d) Q27. (d) Q28. (a) Q29. (d) Q30. (a)
Q31. (c) Q32. (c) Q33. (c) Q34. (a) Q35. (b) Q36. (c) Q37. (c)
Chapter 19
Q1. (a & d) Q2. (a) Q3. (b) Q4. (d) Q5. (b) Q6. (a) Q7. (c) Q8. (a)
Q9. (d) Q10. (d) Q11. (b) Q12. (c) Q13. (a) Q14. (b) Q15. (a) Q16. (b)
Q17. (a) Q18. (b) Q19. (a) Q20. (b) Q21. (d) Q22. (d)
Chapter 20
Q1. (b) Q2. (a) Q3. (a) Q4. (c) Q5. (a) Q6. (b) Q7. (b) Q8. (a)
Q9. (b) Q10. (a) Q11. (d) Q12. (c) Q13. (b) Q14. (b) Q15. (b) Q16. (d)
Chapter 21
Q1. (a) Q2. (b) Q3. (b) Q4. (b) Q5. (b) Q6. (b) Q7. (d) Q8. (b)
Q9. (a) Q10. (c) Q11. (a) Q12. (a) Q13. (a) Q14. (d) Q15. (a) Q16. (b)
Chapter-22:
Q1. (a) Q2. (a) Q3. (d) Q4. (b) Q5. (a) Q6. (a) Q7. (b) Q8. (d)
Q9. (c) Q10. (b) Q11. (b) Q12. (c) Q13. (b) Q14. (d) Q15. (d) Q16. (c)
Q17. (d) Q18. (a) Q19. (a) Q20. (c) Q21. (c) Q22. (c) Q23. (b) Q24. (c)
Q25. (d) Q26. (b) Q27. (c) Q28. (a) Q29. (d) Q30. (a)
Index 809
Index
Symbols Anti-ferromagnetic Materials 688 Beating 270

Antiferromagnetism 686 Bending Moment 800
(Hi-Tc) Superconductivity 726
Anti-node 261 Bending of Beam 799
Antireflecting Coatings 34 Beta Decay 463
A Anti-symmetric Wave Function 454 Betatron 484
Absolute Error 781 Application of Superconductivity B–H Curve 699
Absolute Zero 716 726 Biaxial Crystal 128
Absorbing Antireflecting Coatings Applications of Binding
35 Lasers 169 Energy 455
Absorption and Dispersion of Nanotechnology 768 Energy Per Nucleon 456
Ultrasonic Waves 290 Nuclear Reactor 476 Bioacoustics 295
Absorption Coefficient 299 Optical Fibre Couplers 198 Biopolymers 759
Absorption of Radiation 156 Photoconductivity 675 Biot’s Polariscope 127
Acceptance Angle 190 SEM 224 Biquartz Polarimeter 139
Accuracy 780 X-rays 740 Birefringence 141
Acoustically Good Halls 296 Aragonite 128 Blackbody Radiation 552
Acoustic Environments 297 Arc Discharge Method 757 Body Centered Cubic 520
Acoustics 284 Architectural Acoustics 295 Body Centered Cubic Lattice 524
Acoustics of Buildings 296 AR Coatings 35 Bohr’s
Active System 163 Artificial Radioactivity 460 Magnetons 453, 683
Activity 461 Aspect Ratio 753 First Orbit 599
Addition of Velocities 406 Atomic Bose Einstein
Adequate Loudness 300 Numbers 457 Condensation 724
Agglomeration 763 Nucleus 451 Statistics 454, 615
Air Borne Noise 300 Radius 525 Bosons 615, 724
Allowed Modes 194 Attenuation 195 Bottom-Up Approach 754
Alnico Alloys 701 Attenuation Coefficients 245 Bragg’s Law 531, 561
Alpha Decay 462 Bragg’s X-ray Spectrometer 532,
Amorphous Solids 518 B 564
Ampere’s Circuital Law 338 Ball Milling 760 Brewster’s
Amplitude 234 Bandwidth 187 Angle 125
Analyser 124 Barium 289 Law 125
Angle Resolved Low Coherence Basis 517 Bright Fringes 12
Interferometry 36 BCS Theory 723, 724 Bright Rings 24
Angular Momentum 452 b¯ Decay 464 Brillouin Zones 659
Angular Simple Harmonic Motion b+ Decay 464 Bubble Chamber 469
234 Beams 799 Bucky Balls 765
Angular Width 89 Beat Frequency 270 Bulk Modulus 796
810 Index
C Condensation 763 D
Condition for Maxima 17
Calcite 128 Damped Harmonic Oscillator 243
Condition for Minima 18
Calcite Crystal 128 Damped Motion 245
Conducting Medium 344
Canada Balsam Layer 134 Dark Fringes 13
Conduction Electrons 634
Cantilever 799, 800 Dark Rings 25
Conductors or Metals 663
Capacitor 337 Davisson-Germer Experiment 560
Conservation Laws 471
Carbon Dioxide Gas Laser 167 Dead Time 467
Conservation of Energy 8
Carbon Nanotubes 755 deBroglie
Constancy of Speed of Light 400
Catalyst Free Growth 755 Wave Group 597
Constant Current 226
Certainty 781 Wavelength 655
Constant Height 226
Chain Reaction 474 Waves 559, 560
Construction of the Position 2
Characteristic of Laser Light 161 Debye Length 478
Constructive Interference 2, 7
Characteristics of the Wave Function Degrees of Freedom 751
Continuity Equation 342
602 Del Operator 330
Charge 452 Continuous Beam 799 Dental Care 293
Density 329 Continuous X-ray Spectrum 737 Destructive Interference 3, 7
Independence 455 Controlled Chain Reaction 475 Detection of Ultrasonic Waves 291
Mass and Size 452 Controlled Fusion 477 Dextro-rotatory Substance 138
Chemical Vapour Deposition (CVD) Conventional Photography 171 Diagnosis 292
Technique 758, 766 Cooper Pair 723 Diamagnetic Materials 685, 688
Chirality 757 Cooper Pair Wavefunction 724 Diamagnetism 686
Chromatic Aberration 790 Coordination Number 523 Diamond Structure 523
Circular Aperture 72 Copenhagen Interpretation 551 Dielectric 313
Circularly Polarised Wave 123 Core Electrons 634 Dielectric Constant 313
Cladding 186 Corpuscular 559 Dielectric – Dielectric Boundary
Classical Theory of Corpuscular’ Nature of Light 559 Conditions 340
Diamagnetism 689 Correction of Chromatic Aberrations Dielectric
Ferromagnetism 696 792 Loss 319
Paramagnetism 693 Corresponding Points 85 Medium 344
Clausius-Mosotti Equation 321 Coulomb Effect 459 Polarisation 315
Cleaning 293 Coulomb Gauge 341 Waveguide 188
Coaxial Cable 364 Covalent Bond 530 Differential Equation of SHM 234
Coaxial Capacitor 338 Critical Angle 187 Diffraction and Interference 64
Coercivity 698 Critical Damped Motion 244 Diffraction
Coherence 4 Critical Size of Nucleus 475 Grating 84
Coherence Length 725 Critical Temperature 716 Of Light 63
Coherence Time and Coherence Crystalline Solids 518 Pattern 88
Length 5 Crystallographic Axes 518 Disadvantages of Nanotechnology
Coherence Volume 725 CsCl 523 770
Coherent: 161 CT (Computerized Tomography) Discovery of Neutron 469
Coherent Scattering 561 292 Disintegration Energy 472
Coherent Sources 5 Curie–Weiss Law 695, 697 Disintegration or Decay Constant
Collimated: 161 Curl 332 460
Compositional Defect 537 Current Density 342 Dispersion Relation 354, 358
Compressive Stress 794 CVC 762 Dispersive Power 97
Compton CVD 762 Displacement and Pressure
Effect 561, 562 Cyclotron 482 Amplitude 271
Scattering 562 Cyclotron Frequency 211 Displacement Current 351
Shift 565 Cylindrical Coordinate System 330 Displacement Current Density 348
Index 811
Distributed Load 799 Electronic Ferromagnetic Materials 685, 688,

Divergence 331 Charge 328 689
Division of Defect 537 Fibre Optics Communication 193
Amplitude 16 Polarisation 317 Finite Potential Barrier 610
Wavefront 10 Electronic Speckle Pattern Finite Potential Step 607
Domains 696 Interferometry (ESPI) 36 First London Equation 721
Doppler Effect 264 Electron Lattice Interactions 724 Fixed Beam 799
Doppler Shift 264 Electron Optics 208 Flatness Interferometers 32, 33
Double Refraction 129 Electrostatic Boundary Conditions Flexible Endoscopes 275
Double Slits 80 339 Flux 331
Drift Tubes 482 Electrostatic Focusing 219 Focusing due to Walls and Ceiling
Drift Velocity 214 Electrostriction 289 300
Dushman’s Equation 645 Electromagnetic (EM) Wave Food Technology 294
Dynamic Errors 780 Propagation in Free Space 352 Forbidden Bands 655
Elliptically Polarised Light 135 Force Constant 235
E Elliptically Polarised Wave 123 Forced Vibrations 247
Embossed Hologram 176 Four-Level System 160
Echelon Effect 300
Endoergic Reaction 472 Fractional Refractive Index 193
Echoes 300
Endoscopy 274 Frame of Reference 395
Effective Mass 660
Energy Band 654 Fraunhofer Diffraction 65
Eigen Function 604
Energy Band Diagram 666 Free Electron Gas 634
Eigen Values 604, 606
Energy E versus Wave Number k Frenkel Defect 534, 535
Einstein’s Coefficient of Absorption
Diagram 657 Frequency 234
of Radiation 156
Energy Loss 698 Fresnel Diffraction 65
Einstein’s Coefficient of Spontaneous
Engineering Acoustics 295 Fresnel’s Biprism 10
Emission of Radiation 157
Entropy 718 Fresnel ’s Half-Period Zones 66
Elasticity 794
Errors Occurring in Arithmetic Fresnel Zones 65
Electrical
Operations 783 Fringes 11
Conductivity 635, 636
Ether Drag Hypothesis 400 Fringe Width 22
Properties 759
Excited State 6 Frozen-Fringe Holography 36
Resistivity 717
Exoergic 472 Fullerene 765
Electric
Exoergic Reaction 472 Full Width Half Maximum (FWHM)
Arc Discharge Method 766
External Reflection 9 223
Field 335
Extraneous Noise 300 Fundamental Frequencies 263
Polarization 343
Extraordinary Ray (E-ray) 128 Fundamental Theorem for Gradient
Potential 335
Extrinsic Semiconductor 667 333, 334
Property 453
Extrinsic Sensors 197 Fundamental Theorem of Calculus
Quadrupole Moment 453
333
Susceptibility 316, 321
Electrodeposition 764 F Fundamental Translation Vectors
519
Electromagnetic (EM) Spectrum Face Centered Cubic Lattice 520,
Fundamental Vibration 262
790 524
Electromagnetic Energy Density Falsification 177
360 Far-field Diffraction 65 G
Electromagnetic Wave 121 Fermi-Dirac Distribution 615 Gain Factor 673
Electromagnetic Waveguides 363 Fermi-Dirac Statistics 454, 615 Galilean Acceleration
Electromagnetism 343 Fermi Energy 616, 640 Transformations 397
Electron Capture 464 Fermi Level 640 Galilean Transformation 396
Electron Concentration 664 Fermions 615 Galilean Transformation for Position
Electron-hole Pairs 672 Ferrimagnetic Materials 685 397
812 Index
Galton Whistle Method 289 Hybrid Modes 364 Ionisation Chamber 466
Gamma Decay 465 Hydrogen Bond 531 Iron–Silicon Alloys 701
Gas Condensation Technique 761 Hydrolysis 763 Irradiance 195
Gas Laser 166 Hygiene Safely 293 Isotope Effect 720
Gauss’s Law in Dielectrics 318 Hysteresis 697, 698
Gauss’s Theorem 334 Curve 699 J
Geiger–Mueller Counter 467 Loss 698
Jacket 186
Geodetic Standard Baseline 36
Just Resolved 92
Geometrical Image 75 I
Gradient 331
Grain Orientation 701
Iceland Spar 128 K
Ignition Temperature 479
Grating Element 85 Kronig-Penney Model 655
Image Contrast 222
Green’s Theorem 334 Kundt’s Tube Method 291
Image Formation in SEM 222
Gross Errors 779
Imaging Interferometry 35
Group Velocity 565, 566
Incoherent Scattering 562 L
Guiding Centre 213 Langevin’s Theory 689
Incoherent Source 6
Gyration Velocity 214 Laplace’s Equation 337
Independent Particle Model 458
Gyromagnetic Frequency 211 Laplacian 330
Indoor Acoustics 297
Induced Electric Dipole 315 Larmour Frequency 691
H Induced Radioactivity 460 Larmour Radius 212
Haidinger’s Fringes 30 Inertial Confinement Fusion 479 LASER 156
Half-Life Time 461 Inertial Frame of Reference 396 Laser Ablation Method 757, 766
Half-Period Zones 66 Infinite Potential Well 600, 605 Laser Beat Wave Accelerator 486
Half-Wave Plate 133 Infrasonic Waves 284 Laser Cooling 170
Hall Coefficient 670, 671 Inorganic Nanotube 759 Laser Fusion 480
Hall Effect 669 Insulators 662 Laser Wake Field Accelerator 487
Hall Voltage 670, 671 Intensity 736 Lateral Chromatic Aberration 792
Hamilitonian Operator 604 Intensity of Magnetisation (I) 686 Lateral Shift 14
Hard Ferrites 702 Intensity of Sound 268 Lattice 517
Hard Materials 686 Interatomic Attractive Forces 528 Lattice Constants 519
Harmonic Oscillator 236, 237, 613 Interatomic Repulsive Forces 528 Lattice Parameters 517
Heisenberg Uncertainty Principle Interference Coatings 34 Lattice Planes 521
596 Interference in Thin Films 19 Laue Method 532
Helium-Neon Laser 166 Interference Lithography 37 Laurent Saccharimeter 140
Helix 212 Interference of Light 1 Laurent’s Half-Shade Polarimeter
Hermite Differential Equation 613 Interference of Sound Waves in 138
High Bit Rate 194 Time: Beats 270 Law of Refraction 341
Hole Concentration 665 Interference Pattern 1 Laws of Radioactive Disintegration
Holocameras 178 Interference Principal Maxima 86 460
Holographic Data Storage 178 Internal Reflection 9 Lawson Criterion 479
Holographic Interferometry (HI) Interplanar Spacing 524 Leavo-Rotatory Substance 138
36 Interstitial 537 Length Contraction 403
Holography 170 Intrinsic Semiconductor 664 Lens Aberrations 790
Hooke’s Law 795 Intrinsic Sensors 196 Light Vector 122
Horizontal Oscillations 241 Inverse Lorentz Transformation Limitations of Nanotechnology 770
Huygens-Fresnel Principle 64 Equations 403 Linear
Huygens’ Principle 2 Ion Cores 634 Accelerator 481
Huygens’ Theory of Double Ionic Bond 529 Charge Density 329
Refraction 131 Ionic Polarisation 317 Dispersive Power 98
Index 813
Harmonic Oscillator 600 Magnification 222 Multicoating 35

Polarised Light 135 Malus’ Law 124 Multi-Focus Behaviour of a Zone
Polarised Wave 122 Many Body Effects 717 Plate 70
Simple Harmonic Motion 234 Mass 452 Multimode Fibres 189
Line Spectrum 738 Mass Defect 455 Multimode Graded Index Fibres 189
Liquid-Drop Model 473 Mass Energy Relation 409 Multimode Step Index Fibres 189
Lissajous Figures 272 Mass-String System 241 Multiple Beam Superposition 9
Live Holography Interferometry 36 Matrix Mechanics 552 Multiple-Walled Nanotubes 756
Logarithmic Decrements 246 Matter Waves 559, 560 Multi-Walled Carbon Nanotubes
London Equations 720 Max Planck’s Hypothesis 552 757
London Penetration Depth 722 Maxwell-Boltzmann Statistics 615 Mumetal 701
Longitudinal Chromatic Aberration Maxwell’s Musical Acoustics 295
791 Equations 343
Longitudinal Equations in Conducting N
Coherence 4 Medium 357 NaCl 523
Strain 795 Equations in Isotropic Dielectric Nanomaterials 751
Stress 794 Medium 355 Nanoparticles 760
Wave 2 First Equation 346 Nano Pollution 771
Lorentz Contraction 404 Fourth Equation 347 Nano Robotics 770
Lorentz–Drude Theory 635 Second Equation 346 Nanoscience 752
Lorentz Field 321 Third Equation 347 Nanostructures 751
Lorentz-FitzGerald Contraction Mean Free Path 299 Nanotechnology 752
404 Meissner Effect 717 Nanowire 753
Lorentz FitzGerald Contraction Meson 455 Natural Frequencies 262
Hypothesis 401 Meson Theory 455 Nature of Light 559
Lorentz Force 208 Metallic Bond 530 Nd-YAG Laser 165
Lorentz Transformation 401 Michelson-Morley Experiment 398 Nearest Neighbour Distance 525
Low–Carbon Steel 700 Michelson’s Interferometer 29 Near-Field Diffraction 65
Microscope 94 Neutron Cross-Section 470
M Microscopy 177 Newton’s Rings 23
Machining 293 Microwaves 186 n-hexane Pyrolysis 758
Magnetic Miller Indices 521 Nickel–Iron Alloys 701
Circuits 700 Missing Orders in Diffraction Pattern Nicol Prism 132, 133
Confinement 480 84 Node 261
Field Strength 338 MOCVD 755 Nonaxial Point 74
Flux Density 338 Modes in Waveguides 363 Non-Inertial Frame of Reference
Induction 338 Modified Radiation 570 396
Materials 685 Moified Wavelength 570 Non-Invasive Therapeutic
Moment of an Electron 687 Molecular Beam Epitaxy 755 Applications 292
Polarization 343 Molecular Bond 530 Non-Magnetic Dielectric Medium
Property 453 Molecular Polarisability 321 356
Storage 702 Momentum Space 644 Non-polar Dielectrics 314
Susceptibility 686 Monochromatic 161 Nonreflecting Coatings 34
Vector Potential 341 Monochromatic Light 11 Non-Uniformly Distributed Load
Magnetism 685 Monomode Fibre 194 799
Magnetization 343 Moore’s Law 751 Normalised 601
Magnetostatic Focusing 220 Moseley’s Law 738 Normalised Frequency 194
Magnetostriction Effect 285 MRI (Magnetic Resonance Imaging) Normal Population 157
Magnetostriction Method 285 292 Normal Spectrum 98
814 Index
n-type Extrinsic Semiconductor 667 Particle Accelerators 481 Polarisation 121

n-type Semiconductors 663 Particle in a Box 605 By Double Refraction 129
Nuclear Path Difference 3 By Refraction 128
Energy 474 Penetrating Power 736 By Scattering 134
Fission 473 Penetration Depth 722 Density 315
Forces 454 Percentage Error 782 Of Dielectrics 315
Fusion 476 Permeability 687 Polarised Light 122
Liquid Drop Model 459 Phase 234 Polariser 133
Magic Numbers 458 Difference 3 Polarising Angle 125
Magnetic Moments 453 Reversal Zone Plate 72 Polaroid Filter 124
Radiation Detector 466 Space 644 Population Inversion 157, 159
Reactions 471 Velocity 565 Position Space 644
Reactor 475 Phonons 717 Postulates of Special Theory of
Shell Model 457 Photoconductivity 671 Relativity 401
Spin 453 Photoelasticity 141 Potential Energy Curve 528
Stability 457 Photoelastic Material 141 Powder Method 533
Nucleons 452 Photoelectric Effect 557, 735 Poynting Theorem 361
Numerical Aperture 190 Photon 556 Poynting Vector 361
Physical Acoustics 295 Precision 780
O Physical Significance of Clausius- Primitive Lattice Cell 519
Mosotti Equation 322 Principle of Holography 171
Oblique 799 Probability 601
Piezoelectric
Occupation Index 615 Probability Densities 607
Crystal 285
Occupied Subshells 655 Producing Beats 270
Detection Method 291
Odd-Even Effect 457, 459 Properties of Carbon Nanotubes
Detector 292
Ohm’s Law 635 758
Effect 287
One-Step Hologram 176 Properties of X-Rays 736
Method 287
Operators 604 Proper Time Interval 405
Pion 455
Optical Psychological Acoustics 295
Pitch 212
Coherence Tomography 36 p-type Extrinsic Semiconductor 668
Planck’s
Fibre 186, 187 p-type Semiconductors 663
Oscillators 554
Fibre Connector 197 Pulse Dispersion 195
Quantum Hypothesis 554
Fibre Couplers 197 Pulse Energies 155
Quantum Theory 553
Fibre Sensors 196 Pumping 159, 163
Radiation Law 553
Flatness 32
Plane 799
Properties 764 Q
Diffraction Grating 96
Optic Axis 128
Of Polarisation 123 QDs 767
Orbital Angular Momentum 452
Polarised Light 135 Quality Factor 247
Ordinary Ray (O-ray) 128
Polarised Wave 122 Quantum
Origin of X-Rays 736
Plasma 477 Confinement 752
Oscillatory Motion 233
Plasma–Based Particle Accelerators Dots 767
485 Effects 752
P Plasma Frequency 478 Field Theory 552
Packing Fraction 526 Plasma Wake Field Accelerator 485 Number 606
Pairing of Nucleons 457 PLD 755 Statistics 614
Parallel Plate Capacitor 337 Point Defects 534 Quantum Theory of Free Electrons
Paramagnetic Materials 685, 688 Point Load 799 638
Paramagnetism 686 Poisson’s Equation 337 Quarter-Wave Plate 132
Parity 454 Polar Dielectrics 314 Quartz 128, 287
Index 815
R Schottky Defects 534 -Intensity Level 269

Schrödinger Equation 656 Speed 267
Radioactivity 460
Scintillation Counter 468 Wave and its Velocity 260
Radio Communications 187
Scratch Resistant Coating 769 Waves 284
Radius of Curvature of Plano Convex
Secondary Maxima 87 Sources of Errors 780
Lens 25
Secondary Wavelets 2 Space and Time 2
Rainbow Hologram 176
Second London Equation 722 Spatial Coherence 5
Random Errors 780
Self–Cleaning Glass 768 Spatial Interval 5
Rayleigh Criterion for Resolution 92
Self–Phase Modulation LWFA 487 Specific Charge of an Electron 208
Rayleigh-Jeans Law 553, 555, 556
SEM 755 Specific Heat 718
Recombination Coefficient 672
SEM Components 221 Specific Rotation 138
Reconstruction of Image 172
Semiconductor Detector 469 Spherical Capacitor 338
Red Shift 264
Semiconductor Laser 168 Spherical Polar Coordinate System
Reflection Hologram 175
Semiconductors 662 330
Refractive Index 356
Semi-Reverberant 297 Spin Angular Momentum 452
Refractive Index of a Liquid 25
SEM Principle 221 Spontaneous Emission 156
Relation between Elastic Constants
Sensitive Flame Method 291 Sputtering 762
796
Sensitive Method 291 Stability Curve 457
Relative
Sensitivity 673 Standard Deviation 781
Error 782
Shearing Force 800 Standing Waves 261
Permeability 686
Shearing Stress 795 Standing Waves in Air Columns
Refractive Index 192
Shear Modulus 795 262
Relaxation Time 246
Significance of Maxwell’s Equations Static Error 779
Resolution 222, 780
350 Statistical Property 454
Resonance 286, 300
Simple Cubic 520 Stokes’ Theorem 335
Resonance Condition 262, 486
Simple Cubic Lattice 524 Stopping Potential 558
Resonant Cavity 163
Simple Harmonic Motion (SHM) STPs 755
Response Time 675
234 Strain 794
Retentivity 698
Simple Pendulum 239 Strength 758
Reverberation 296, 300
Simply Supported Beam 799 Stress 794
Reverberation Time 296
Single Stress and Strain 794
Richardson-Dushman Equation 645
Mode Fibre 194 Structure Borne Noise 300
Richardson-Laue-Dushman Equation
Mode Step Index Fibre 188 Superconductivity 716
645
Optical Fibre 186 Superposition of Two SHMs 260
Richardson’s Equation 644
Slit 76 Superposition Principle 336
Rigid Endoscopes 275
-Walled Carbon Nanotubes 756 Supersonic and Shock Waves 266
Rochelle Salt 287
-Walled Nanotubes 756 Surface Charge Density 329
Ruby Laser: Solid State Laser 163
Size 452 Surface Energy 459
Skin Depth 357, 359 Surface Plasmon Resonance Imaging
S Skip Distance 192 765
Sabine’s Formula 298 Smart Materials 769 Surface Plasmons 764
Saccharimeter 140 Soft Ferrites 702 Surgery 292
Scalar Potential 341 Soft Materials 686 Sustained Interference 8
Scanning Electron Microscope Sol-Gel Techniques 763 Swarms 770
(SEM) 220 Sound Symmetric Wave Function 454
Scanning Tunneling Microscope Absorbing Materials 301 Symmetry Effect 459
(STM) 225 Barrier 266 Synthesis of Fullerenes 766
Schemes for Population Inversion Displacement 261 Synthesis of Nanowires 753
159 Insulation 300 Systematic Errors 779
816 Index
T Transverse Nature of Velocity of Sound in Liquid 294

Electromagnetic Waves 355 Vertical Oscillations 242
Telescope 93
Traps 674 Vibration 284
TEM 755
TV Antenna 121 Virtual Coherent Sources 14
Temporal Coherence 4
Two-Level System 159 Volume
Thermal Conductivity 637, 719, 635
Two Step Hologram 176 Charge Density 329
Thermal Properties 718, 759
Twyman-Green Interferometer 33 Energy 459
Thermionic Emission 644
Type-II (Hard) Superconductor 719 Hologram 177
Thick Holograms 177
Type-I (Soft) Superconductor 719 Strain 795
Thickness of Thin Transparent Sheet
15 Types of
Thin Films 16 Acoustics 295 W
Thomson’s Method 209 Crystals 519
Wave 2
Three-Level System 160 Dielectrics 314
Wavefront 2
Threshold Diffraction 64
Wave Function 601, 603, 655
Energy 473 Endoscopes 274
Waveguide 363
Frequency 558 Polarisation 317
Wavelength of Light 13, 90
Time Wave Mechanics 552
Average Holographic U Wave Nature of Light 63
Interferometry 177 Ultrasonic Wave Particle Duality 557
Dependent Wave Equation 604 Hologram 177 Wedge Shaped Film 20
Dilation 405 Transducer 288 Welding 293
Independent Schrödinger Waves 284, 285 White Light Hologram 176
Equation 603 Uncontrolled Chain Reaction 474 Wiedemann–Franz Law 635, 638
Period 234 Uniaxial Crystal 128 Wien’s Law 555, 556
Titanate 289 Uniform Magnetic Field 210 Wigner–Seitz Primitive Cell 519
Topaz 128 Unit Cells 517, 518 Wilson’s Cloud Chamber 468
Top-Down Approach 754 Unmodified Radiation 570 Work Function 558, 644
Tourmaline 128 Unmodified Wavelength 570
Trajectories 210 Unoccupied Subshells 655
Transference 2 X
Unpolarised Light 122
Translational Motion 233 X-Rays 735
Uses of Coaxial Cables 365
Translation Vector 520 X-Ray Spectra 737
Transmission Hologram 174
Transmission Loss 297 V Y
Transmission of Sound 297 Vacancies 534
Young’s Double Slit Experiment 2,
Transmitted Light 26 Vacuum Deposition 762
64
Transverse Electric (TE) Modes Van der Waals Bond 530
Young’s Modulus 795
363 Vaporization 762
Transverse Electromagnetic (TEM) Vapour Liquid Solid 754
Modes 364 Vapour Solid (VS) Method 754 Z
Transverse Magnetic (TM) Modes Variation of Mass with Velocity Zone Plate 68
363 407 Zone Plate and Convex Lens 71

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Engineering

McGraw Hill Education (India) Private Limited

McGraw Hill Education Ofﬁces

Published by McGraw Hill Education (India) Private Limited

Managing Director: Kaushik Bellani

Head––Production: Satinder S Baveja

Visit us at: www.mheducation.co.in

22. Nanoscience and Nanotechnology 751

2.7 Fraunhofer Diffraction by a Single Slit 76

4.9 Holography Versus Conventional Photography 171

7.6 Mass-string System 241

Objective Type Questions 308

11.25 Maxwell’s Equations in Isotropic Dielectric Medium: EM Wave Propogation 355

13.12 Nuclear Fusion 476

Solved Examples 570

18.10 Effect of Traps 674

Objective Type Questions 732

Professor Prem Vrat

Preface to the Second Edition

1.1 YOung’s DOubLe sLit experiment LO1

1.2 COnCept Of Waves anD HuYgens’ prinCipLe LO2

In order to demonstrate the Huygens’ principle, we consider P1 P2 P P1

1.3 pHase DifferenCe anD patH DifferenCe LO3

1.3.3 Relation between Path Difference and

1.4 COHerenCe LO4

1.4.3 Coherence Time and Coherence Length

1.5 COHerent sOurCes LO4

1.5.1 Production of Coherent Light from Incoherent Sources

Pinhole Wavelength Filter

1.6 anaLYtiCaL treatment Of interferenCe LO5

According to the principle of superposition, the resultant displacement (yR) is given by

1.6.1 Condition for Constructive Interference

1.6.2 Condition for Destructive Interference

If a1 π a2, then Imin π 0

I min < a12 + a22

–7p –6p –5p –4p –3p –2p –p 0 p 2p 3p 4p 5p 6p 7p

Therefore, average intensity (Iav) will be obtained as

1.7 COnDitiOns fOr sustaineD interferenCe LO5

1.7.1 Condition of Relative Phase Shift

1.8 muLtipLe beam superpOsitiOn LO6

1.9 interferenCe bY DivisiOn Of WavefrOnt LO6

1.9.1.3 Determination of Distance between Two Virtual Coherent Sources

Then from Eqs. (i) and (ii), we get

1.10 interferenCe bY DivisiOn Of ampLituDe LO6

1.10.1 Interference Due to Thin Films

The path difference between the reflected rays

From Eqs. (xvii) and (xviii), we get

Eye Extended Source Eye

The path difference between the two reflected rays

In right angled DBND,

Condition for Maxima: For constructive interference

1.10.2.3 Determination of Refractive Index of a Liquid

1.10.2.4 Newton’s Rings in Transmitted Light

D = 2t (ii) Figure 1.22

1.10.2.5 Newton’s Rings formed between Two Curved Surfaces

The thickness of the air film at P is Figure 1.24

For nth bright fringe

(a) (b) (c) (d)

On subtracting Eq. (i) from Eq. (ii), we get